Oscillations and Waves

PH-1007 (Physics)
Dr. Gorky Shaw

1.1 Simple harmonic motion

1.1.1 Introduction
1. Periodic motion
Repeated motion along a definite path at regular intervals of time is called
periodic motion. The interval of time is called time interval.

2. Oscillation or vibratory motion

To and fro or back and forth motion repeatedly about a mean position is
called oscillation or vibratory motion.
Such periodic and bounded motion is confined within well-defined limits
on either side of the mean position. The body performing oscillation is
called oscillator.
To represent this type of motion mathematically, sine and cosine functions
are best suited, being both simple periodic and bounded.

3. Simple harmonic motion (SHM)

Oscillatory motion is in which the amplitudes are equal on both sides of
the mean position is called simple harmonic motion (SHM).
Examples of SHM are motion of a spring and simple pendulum.

Figure 1.1: Graphical representation of simple harmonic motion.

 1.1 Simple harmonic motion
4. Necessary condition for SHM
A necessary condition for SHM is a restoring force acting on the body.

F ∝ −y
or, F = −ky (1.1)

where y is the displacement of the body from its mean position. k is

called force constant or spring constant.

1.1.2 Equation of motion and solution

1. Differential equation for SHM.
For any body of mass m, being acted upon by a force F and having
velocity v and acceleration a, the equation of motion is,

F = ma (1.2)

Now,

dv d2 y
a= = 2 (1.3)
dt dt

Therefore, from (1.1) and (1.2) and (1.3), we get

d2 y
m 2 = −ky
dt
d2 y k
=⇒ 2 + y = 0
dt m
2
dy
=⇒ 2 + ω 2 y = 0 (1.4)
dt

k
where ω 2 = is a positive constant.
m

(1.4) is the differential equation representing SHM.

A special model for SHM, the simple pendulum, is discussed later in

Section 1.6.1.
The next step is to find a solution of this equation to characterize the
motion.

2
 1.1 Simple harmonic motion

2. Solution of the equation and expression for the displacement in

SHM
By making an educated guess, we can see that two possible solutions of
(1.4) are

y = eiωt and y = e−iωt (1.5)

which we call y1 and y2 , respectively.

In such cases, it is known that, since y1 and y2 are solutions of (1.4), their
linear combination

y = C 1 y 1 + C 2 y2
or, y = C1 eiωt + C2 e−iωt (1.6)

where C1 and C2 are constants (usually complex numbers) to be deter-

mined from initial conditions, is also a solution of (1.4).

Exercise: Substitute y1 = eiωt , y2 = e−iωt , and y = C1 y1 + C2 y2 in (1.4)

and check that they all are indeed solutions of this equation.

Now, we know that

eiωt = cos ωt + i sin ωt

and, e−iωt = cos ωt − i sin ωt (1.7)

Substituting (1.7) in (1.6), we get

y = C1 (cos ωt + i sin ωt) + C2 (cos ωt − i sin ωt)

=⇒ y = (C1 + C2 ) cos ωt + i(C1 − C2 ) sin ωt (1.8)

We make the replacement of variables:

A1 = C1 + C2
and, A2 = i(C1 − C2 ) (1.9)

Thus, from (1.6), (1.8), and (1.9), General solution of (1.4) is,

y = A1 cos ωt + A2 sin ωt (1.10)

Where A1 and A2 are constants to be determined from initial conditions.

3
 1.1 Simple harmonic motion

For example, we consider the following initial condition: at t = 0, y = 0.

Putting this in (1.10), we get

0 = A1 + 0 =⇒ A1 = 0 (1.11)

y = A2 sin ωt for this initial condition. (1.12)

Exercise: Check that y given in (1.12) satisfies (1.4).

We introduce new variables A and φ, related to A1 and A2 as:

A1 = A sin φ
and, A2 = A cos φ (1.13)

Then, substituting (1.13) in (1.10), we get

y = A sin φ cos ωt + A cos φ sin ωt

or, y = A(sin φ cos ωt + cos φ sin ωt)
or, y = A sin(ωt + φ) (1.14)

Hence, we may rewrite the general solution in a more convenient form,

y = A sin(ωt + φ) (1.15)

where the term (ωt + φ) is called the phase of the SHM, and φ is called
the initial phase or phase constant.

Using the above initial condition: at t = 0, y = 0 in (1.15) we get,

0 = A sin φ =⇒ φ = nπ, n = 0, 1, 2, ... (1.16)

Without loss of generality we can take the value φ = 0.

Therefore, y = A sin ωt for this initial condition. (1.17)

(1.6), (1.10), and (1.15) are all valid forms of the general solutions of
(1.4). However, (1.15) is the most convenient form to represent and study
simple harmonic motion.

4
 1.1 Simple harmonic motion

3. Parameters associated with SHM

The parameters associated with SHM are:
ˆ A → Amplitude of oscillation (maximum displacement from the
mean position).
ˆ ω → angular frequency.
ˆ f → frequency. ω = 2πf
1 2π
ˆ T → Time period of oscillation. T = = .
f ω
4. Expressions for velocity and acceleration associated with SHM
Using the expression for y in (1.15) the velocity v is given by,
dy
v= = Aω cos(ωt + φ) (1.18)
dt q
= ±Aω 1 − sin2 (ωt + φ)
q
= ±ω A2 − A2 sin2 (ωt + φ)
p
= ±ω A2 − y 2 (1.19)

Velocity is represented by either (1.18) or (1.19).

Acceleration a is given by,

d2 y
a = 2 = Aω 2 sin(ωt + φ) (1.20)
dt
= −ω 2 y (1.21)

Acceleration is represented by either (1.20) or (1.21).

1.1.3 Expression for energy in SHM

In order to determine the energy of the oscillator, we first evaluate its kinetic
energy and potential energy separately and then add them to get the total
energy.
Kinetic energy of the oscillator,
1 1
K = mv 2 = mω 2 (A2 − y 2 ) (using (1.19)) (1.22)
2 2
Clearly, Maximum and minimum values of kinetic energy are at the mean
position and extreme position, respectively.

5
 1.1 Simple harmonic motion

1
Kmax = mω 2 A2 at y = 0
2
Kmin = 0 at y = A (1.23)

Potential energy at displacement y is the work done against the restoring

force in moving the body from the mean position to this position:
Z y Z y
P = −F dy = kydy
0 0
1
= ky 2
2
1
= mω 2 y 2 (since k = mω 2 ) (1.24)
2
Clearly, Maximum and minimum values of potential energy are at the extreme
position and mean position, respectively.

1
Pmax = mω 2 A2 at y = A
2
Pmin = 0 at y = 0 (1.25)

The total (mechanical) energy,

1
E = K + P = mω 2 A2 (1.26)
2
is conserved (constant).

6
 1.2 Damped Harmonic Oscillations

1.2 Damped Harmonic Oscillations

As discussed in Section 1.1, the amplitude of oscillations remains constant in
case of simple harmonic oscillations (SHM).
When the amplitude of oscillations of an oscillator reduces due to an external
force, the oscillator and its motion are said to be damped.

1.2.1 Equation of motion

The forces acting on such an oscillator are:

1. Restoring force:

FR ∝ −y
or, FR = −ky (1.27)

2. Damping force: A retarding (decelerating) force that is proportional to

its velocity v.
dy
FD ∝ −
dt
dy
or, FD = −b (1.28)
dt

Hence, the equation of motion for a oscillator of mass m is

d2 y
m = F R + FD
dt2
d2 y dy
or, m 2 = −ky − b
dt dt
d2 y dy
=⇒ m 2 +b + ky = 0
dt dt
d2 y b dy k
=⇒ + + y=0 (1.29)
dt2 m dt m
b k
Let = 2r and = ω02 . Then we get
m m

d2 y dy
2
+ 2r + ω02 y = 0 (1.30)
dt dt
This is the equation of motion of a damped harmonic oscillator.
r is called the damping constant, and ω0 is called the natural frequency of the

7
 1.2 Damped Harmonic Oscillations

oscillator. Natural frequency is the frequency of the oscillator in the absence of

any damping (i.e., for r = 0).

1.2.2 Solution of the equation and expression for the displacement

In Section 1.1.2, we noticed that the solution of the equation for SHM involves
exponentials. Similarly, for the case of damped motion, we consider a possible
solution of the form

y = eαt (1.31)

where α is a parameter of motion and is to be determined.

(Note that in case of SHM we directly guessed the solutions y1 = eiωt and
y2 = e−iωt , which worked. However, in the case of damped motion, we consider
a more general exponent α which might depend both on r and ω0 .)
Substituting (1.31) in (1.30), we get

(α2 + 2rα + ω02 )eαt = 0 (1.32)

Thus, we get

α2 + 2rα + ω02 = 0
q
=⇒ α = −r ± r2 − ω02 (1.33)

Hence, we have two possible values of α:

q
α1 = −r + r2 − ω02
q
and α2 = −r − r2 − ω02 (1.34)

and hence two possible expressions for y:

√
r2 −ω02 )t
y1 = e(−r+
√
(−r− r2 −ω02 )t
and y2 = e (1.35)

8
 1.2 Damped Harmonic Oscillations

Hence, the general solution of (1.30) is given by

y = A1 y1 + A2 y2
√ √
(−r+ r2 −ω02 )t (−r− r2 −ω02 )t
or y = A1 e + A2 e (1.36)

where A1 and A2 are constants of motion, to be determined from initial condi-

tions.

Based on the extent of damping, there are three distinct cases of damped
harmonic motion:
1. r2 < ω02 : underdamped motion.
2. r2 > ω02 : overdamped motion.
3. r2 = ω02 : critically damped motion.
As we will see subsequently, in the different cases, the displacement varies
with time as represented graphically in Figure 1.2.

Figure 1.2: Displacement versus time in damped motion (Source: math24.net).

1. Underdamped motion (or light damping)

In this case, r2 < ω02 , which means (r2 − ω02 ) < 0. We introduce a new
parameter ω given by
q
2
q
2
√ q 2
r2 − ω0 = −(ω0 − r2 ) = −1 ω0 − r2 = iω (1.37)

9
 1.2 Damped Harmonic Oscillations
p
where ω = ω02 − r2 is a real number and represent the angular frequency
of damped oscillations. Clearly, ω < ω0 .
Hence, we can rewrite (1.36) as

y = A1 e(−r+iω)t + A2 e(−r−iω)t
y = A1 e−rt eiωt + A2 e−rt e−iωt
= e−rt (A1 eiωt + A2 e−iωt )
= e−rt [A1 (cos ωt + i sin ωt) + A2 (cos ωt − i sin ωt)]
= e−rt [(A1 + A2 ) cos ωt + i(A1 − A2 ) sin ωt] (1.38)

Now, since we are discussing oscillations, which are very real, the displace-
ment y is certainly a real quantity. Hence, both (A1 + A2 ) and i(A1 − A2 )
must be real quantities.
Accordingly, we make the replacement of variables:

A1 + A2 = A0 sin φ
and i(A1 − A2 ) = A0 cos φ (1.39)

Then, from (1.38) and (1.39), we have the general solution for under-
damped harmonic oscillations

y = A0 e−rt sin(ωt + φ) (1.40)

Figure 1.3: Graphical representation of underdamped oscillations.

10
 1.2 Damped Harmonic Oscillations

2π 2π
which represents harmonic oscillations with time period T = =p 2
ω ω0 − r 2
As mentioned earlier, ω0 is the natural frequency (without any damping)
and ω is the damped frequency.
The amplitude of oscillations

A = A0 e−rt (1.41)

decays exponentially with time.

2. Overdamped motion (or hard damping)

p
In this case r2 > ω02 . Hence, r2 − ω02 is a real number.
p
Let r2 − ω02 = β.

Then for the overdamped case, the general solution given by (1.36) can
be rewritten as

y = A1 e(−r+β)t + A2 e(−r−β)t (1.42)

Now clearly, r > β.

Therefore, both the exponential terms in (1.42) decay with time. Thus,
in this case there are no oscillatory terms, and the resultant motion is
non-oscillatory. This type of motion is also called aperiodic motion.

3. Critically damped motion

Critical damping occurs when the damping constant is equal to the nat-
ural frequency of the oscillator. In this case r = ω0 .
Then the general solution given by (1.36) is reduced to

y = A1 e−rt + A2 e−rt = (A1 + A2 )e−rt = Ce−rt (1.43)

where we introduce a new parameter C = A1 + A2 , which is a single con-

stant term in the expression of y. Hence, (1.43) is insufficient as a general
solution of a second order equation (which must include two distinct con-
stants, as seen in all the earlier general solutions).
p
Therefore, in this case, we instead consider r2 − ω02 = h → 0 but 6= 0.

11
 1.3 Energy decay in underdamped harmonic oscillations

Then, we have

y = A1 e(−r+h)t + A2 e(−r−h)t
= e−rt (A1 eht + A2 e−ht )
= e−rt [A1 (1 + ht + ...) + A2 (1 − ht + ...)]
= e−rt [(A1 + A2 ) + h(A1 − A2 )t] (1.44)

In the above expression, we have replaced the exponentials involving h in

the power with the respective series expansions, and ignored the higher
order terms since h is very small.
We further rewrite (1.44) as

y = e−rt (P + Qt) (1.45)

which is the general solution for critically damped motion.

where P = A1 + A2 and Q = h(A1 − A2 ). We can find P and Q from

initial conditions.
Initially, y given by (1.45) increases due to the linear term in t. But
with increasing t, the exponential decay term dominates, and y decreases
rapidly.

Critical damping provides the quickest approach to zero amplitude for a

damped oscillator. With overdamping, the approach to zero is slower. In
case of underdamping, zero is reached more quickly, but is crossed, and there
is oscillation about zero.

1.3 Energy decay in underdamped harmonic oscillations

For underdamped oscillations, we have (from (1.40))

y = A0 e−rt sin(ωt + φ) (1.46)

Therefore,
dy
v= = ωA0 e−rt cos(ωt + φ) − rA0 e−rt sin(ωt + φ)
dt
= A0 e−rt [ω cos(ωt + φ) − r sin(ωt + φ)] (1.47)

Now, for small damping, r << ω.

12
 1.4 Parameters associated with underdamped harmonic oscillations

Therefore, in this case,

dy
v= ≈ A0 e−rt [ω cos(ωt + φ)]
dt
= ωA0 e−rt cos(ωt + φ) (1.48)

Therefore, kinetic energy of the oscillator,

1 1
K = mv 2 = mω 2 A20 e−2rt cos2 (ωt + φ) (1.49)
2 2
and its potential energy
1 1
P = ky 2 = mω 2 y 2
2 2
1
= mω 2 A20 e−2rt sin2 (ωt + φ) (1.50)
2
Therefore, total energy
1
E = K + P = mω 2 A20 e−2rt [cos2 (ωt + φ) + sin2 (ωt + φ)]
2
1
= mω 2 A20 e−2rt (1.51)
2
Hence, we may write

E = E0 e−2rt (1.52)

1
where E0 = mω 2 A20 us the energy at t = 0.
2
Thus, the energy of a damped harmonic oscillator decays exponentially with
time.
Note that while amplitude decay is proportional to e−rt , energy decay is
proportional to e−2rt . This is expected because the energy of oscillation is
proportional to the square of the amplitude, as seen in the above relations.

1.4 Parameters associated with underdamped harmonic

oscillations
1. Logarithmic decrement (λ): it is ratio between successive amplitudes,
at time intervals of T /2.
The displacement for underdamped oscillations is given by (1.40):

y = A0 e−rt sin(ωt + φ) = A sin(ωt + φ) (1.53)

13
 1.4 Parameters associated with underdamped harmonic oscillations

where A0 initial amplitude at t = 0, and A = A0 e−rt is the amplitude at

time t.

Figure 1.4: Displacement versus time for underdamped oscillations. Note that
successive amplitudes are separated by time intervals of T /2. Suc-
cessive amplitudes of the same sign are separated by time intervals
of T .

We consider the time tn and tn+1 such that

tn+1 = tn + T /2 (1.54)

Then the respective amplitudes at tn and tn+1 are

An = A0 e−rtn
and An+1 = A0 e−rtn+1 = A0 e−r(tn +T /2) (1.55)

Therefore,

An A0 e−rtn
= −r(t +T /2)
= erT /2 = constant (1.56)
An+1 A0 e n

erT /2 = d is a constant of motion called decrement.

rT
loge d = =λ (1.57)
2
is called the logarithmic decrement.
Alternative definition of λ: Instead of considering successive ampli-
tudes (either positive or negative), only positive (or only negative) am-
plitudes may be considered. In this case, successive positive (or negative)

14
 1.4 Parameters associated with underdamped harmonic oscillations

amplitudes are separated by a time difference of T , instead of T /2. Ac-

cording to this definition,

d = erT and λ = loge d = rT (1.58)

2. Relaxation time for amplitude (τA ): it is the time when the amplitude
of oscillations decreases to 1/e times its initial value.
1
If we put t = in the expression for A, then we get
r

A0
A = A0 e−r(1/r) = A0 e−1 = (1.59)
e

1
Thus, τA = is the time in which the amplitude reduces to the initial
r
amplitude, and hence is the relaxation time.
In terms of the relaxation time τA ,

A = A0 e−t/τA (1.60)

3. Relaxation time for energy (τ ): it is the time when the energy of

oscillations decreases to 1/e times its initial value.
From (1.52),

E = E0 e−2rt (1.61)

1
Therefore at t = ,
2r
E0
E = E0 e−2r1/2r = E0 e−1 = (1.62)
e

1
Thus, τ = is the relaxation time for energy of oscillations.
2r
In terms of the relaxation time τ ,

E = E0 e−t/τ (1.63)

4. Power dissipation (P ): it is the rate of decrease of energy due to

15
 1.5 Forced oscillations

damping.
dE d 1 E
P =− = − (E0 e−t/τ ) = (E0 e−t/τ ) =
dt dt τ τ
E
Therefore, P = (1.64)
τ

5. Quality factor (Q): quality factor of an oscillator expresses its efficiency.

Energy stored per cycle
Q = 2π
Energy dissipated per cycle
E
= 2π where T is the time period of oscillations
PT
2π E
= ·
T P
= ωτ (1.65)

Thus, Q = ωτ .
For an ideal oscillator, τ → ∞ (no loss of energy). Therefore, in this
Q → ∞.

1.5 Forced oscillations

Objects can be forced to oscillate, most easily at their natural frequency.
Natural frequency is the frequency at which a system would oscillate if there
were no driving or damping forces (i.e., it is unaffected by external forces).
Applying a periodic driving force on a simple harmonic oscillator puts energy
into the system at the driving frequency.
When a body being acted upon by an external periodic force starts to oscillate
with the frequency of the external periodic force, the body is said to undergo
forced oscillations.

1.5.1 Equation of motion

Let an oscillator be subjected to an external periodic force, in addition to
restoring force and damping force. The forces acting on such an oscillator are:

1. Restoring force:

FR ∝ −y
or, FR = −ky (1.66)

16
 1.5 Forced oscillations

2. Damping force: A retarding (decelerating) force that is proportional to

its velocity v.
dy
FD ∝ −
dt
dy
or, FD = −b (1.67)
dt

3. External periodic force:

Fext = F0 sin ωt (1.68)

where F0 and ω are the amplitude and (angular) frequency, respectively,

of the periodic driving force.

Hence, the equation of motion for a oscillator of mass m is

d2 y
m 2 = Fext + FR + FD
dt
d2 y dy
or, m 2 = F0 sin ωt − ky − b
dt dt
d2 y dy
=⇒ m 2 +b + ky = F0 sin ωt
dt dt
d2 y b dy k F0
=⇒ + + y= sin ωt (1.69)
dt2 m dt m m
b k F0
Let = 2r, = ω02 and = f0 . Then we get
m m m

d2 y dy
2
+ 2r + ω02 y = f0 sin ωt (1.70)
dt dt
This is the equation of motion for forced oscillations.
r is the damping constant, and ω0 is natural frequency of the oscillator.
f0 is the driving acceleration.

1.5.2 Solution of the equation and expression for the displacement

The complete solution of the equation for forced oscillations (1.70) consists of
two parts: a transient solution, and a steady-state solution.

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 1.5 Forced oscillations

1. Transient solution
It is the solution of the equation obtained by setting the RHS of (1.70)
to zero, i.e.,

d2 y dy
2
+ 2r + ω02 y = 0 (1.71)
dt dt
which is identical to (1.30) and represents damped harmonic motion. As
we have discussed in Section 1.2, the displacement represented by this
equation decays exponentially with time, and hence is not of interest when
considering forced oscillations over relatively long durations of time.

2. Steady-state solution Once the transient displacements given by (1.71)

have died down, we have the steady-state motion of the system under
forced oscillations. The steady-state solution is to be obtained by trial
method.
Physically, we would expect that the driven system will oscillate with the
frequency of the driving force, and there will be a phase difference between
the driving force and the driven system.
We thus consider a solution of the form

y = A sin(ωt − θ) (1.72)

which satisfies both the conditions. It represents oscillations with (angu-

lar) frequency ω, same as that of the driving force, with a phase difference
of θ. Here A is the amplitude of the forces oscillations and θ is the phase
difference between the driving force and the driven system.
From (1.72),
dy
= Aω cos(ωt − θ)
dt
d2 y
and 2
= −Aω 2 sin(ωt − θ) (1.73)
dt

Using (1.72) and (1.73) in (1.70), we get

−Aω 2 sin(ωt − θ) + 2rAω cos(ωt − θ) + ω02 A sin(ωt − θ)

= f0 sin ωt
= f0 sin[θ + (ωt − θ)]
= f0 sin θ cos(ωt − θ) + f0 cos θ sin(ωt − θ) (1.74)

Now, (1.74) has to be valid for all values of t. Hence, the coefficients of

18
 1.5 Forced oscillations

sin(ωt − θ) and cos(ωt − θ) on the LHS and RHS must be separately equal
to each other. That is,

−Aω 2 + ω02 A = f0 cos θ (1.75)

and

2rAω = f0 sin θ (1.76)

Squaring and adding (1.75) and (1.76), we get

A2 (ω02 − ω 2 )2 + 4r2 ω 2 = f02

 
(1.77)

Therefore, from (1.77), the amplitude of driven oscillations in the steady-

state is given by
f0
A= p 2 (1.78)
(ω0 − ω 2 )2 + 4r2 ω 2

And hence, from (1.72), the steady-state solution of (1.70) is

f0
y=p 2 sin(ωt − θ) (1.79)
(ω0 − ω 2 )2 + 4r2 ω 2

The phase difference (θ) between the driving force and the driven system
may be obtained by taking the ratio of (1.75) and (1.76):
2rAω f0 sin θ
2 =
−Aω 2 + ω0 A f0 cos θ
 
−1 2rAω
=⇒ θ = tan (1.80)
−Aω 2 + ω02 A

When the driving frequency matches the natural frequency of the

driven system, that is when ω = ω0 , we have
 
2rAω π
θ = tan−1 = tan−1 (∞) = (1.81)
0 2

1.5.3 Conditions on the amplitude A

1. ω << ω0 (low frequency of driving force)

19
 1.5 Forced oscillations

In this case, from (1.78),

f0 F0 /m F0
A≈ 2 = = (1.82)
ω0 k/m k
Thus, in this case, the amplitude of oscillations depends only on the restor-
ing force constant k.

2. ω >> ω0 (high frequency of driving force)

In this case, from (1.78),
f0
A≈ (considering small r << ω, i.e., low damping)
ω2
F0 /m
=
ω2
F0
= (1.83)
mω 2
Thus, the amplitude decreases with increasing mass and drive frequency.

3. ω = ω0 (amplitude resonance) The maximum amplitude of oscillations

occurs when the driving frequency matches the natural frequency of the
driven system. When the two exactly match, that is when ω = ω0 , from
(1.78) we have
f0
A= (1.84)
2rω0
Thus, in this case, amplitude is high for small r (low damping). For
r → 0, A → ∞.
Apparently this corresponds to the maximum possible amplitude Amax .
But in presence of damping, Amax does not occur exactly at ω = ω0 .
For Amax , the condition to be satisfied is
dA
=0 (1.85)
dω

20
 1.5 Forced oscillations

From (1.78), it is evident that a simpler condition for the same is

d  2
(ω0 − ω 2 )2 + 4r2 ω 2 = 0

dω
=⇒ 2(ω02 − ω 2 )(−2ω) + 4r2 (2ω) = 0
=⇒ 4ω(−ω02 + ω 2 + 2r2 ) = 0
=⇒ either ω = 0 or (−ω02 + ω 2 + 2r2 ) = 0
(1.86)

ω = 0 is ignored because it implies there is no driving force. Thus, we

have,

(−ω02 + ω 2 + 2r2 ) = 0
=⇒ ω 2 = ω02 − 2r2
q
=⇒ ω = ω02 − 2r2 (1.87)

Therefore, the maximum amplitude Amax occurs at

q
ω = ω02 − 2r2 < ω0 (1.88)

in presence of non-zero damping.

For very low damping (very small r << ω0 ), Amax occurs at

ω ' ω0 (1.89)

This is the case of amplitude resonance.

1.5.4 Sharpness of resonance and Quality factor

1. Resonance When the frequency of the driving force approaches the nat-
ural frequency of the driven system, the amplitude of oscillations becomes
very large. A small-amplitude driving force can produce a large-amplitude
response. This phenomenon is known as resonance.
Examples of resonance are found in tuning forks, pushing a swing, soldiers
marching, NMR (nuclear magnetic resonance) etc.

2. Sharpness of resonance Amplitude at resonance is large for low damp-

ing and small for heavy damping. Further, for low damping, amplitude
drops rapidly as ω shifts away from ω0 . Thus the resonance is sharp. For
heavy damping, the decay of amplitude with ω is more gradual. Thus the

21
 1.5 Forced oscillations

resonance is flat.
The variation of peak amplitude with damping is schematically repre-
sented in Figure 1.5. The figure also indicates the shift in the resonance
frequency farther away from ω0 with increasing damping, as expected
from (1.88).
Increasing damping also has the effect of shifting the phase difference θ
at resonance from 900 to lower values.

Figure 1.5: (a) Effect of damping on amplitude of forced oscillations (Source:

xmDemo). (b) FWHM and Q-factor (Source: researchgate).

22
 1.6 Coupled oscillations

3. Quality factor (Q-factor) The spread ∆ω of the angular frequency at

half the maximum amplitude is known as full width at half maximum or
FWHM or bandwidth.
Quality factor or Q-factor is the defined as the ratio
ω0
Q= (1.90)
∆ω
Q-factor is a measure of the sharpness of resonance.
∆ω small =⇒ Q large =⇒ sharp resonance.
∆ω large =⇒ Q small =⇒ flat resonance.

4. Bandwidth and damping It can be shown that

1
∆ω = 2r = (1.91)
τ
1
where τ = is defined as the relaxation time of oscillations, analogous
2r
to the case of damped harmonic oscillations.

1.6 Coupled oscillations

1.6.1 Revisiting SHM (for small angles) - the simple pendulum
A simple pendulum consists of a spherical bob (considered a point mass) at-
tached to a string (considered massless). It is an idealized mathematical model
of a real pendulum. We consider a pendulum of mass m and length L, as shown
in Figure 1.6. The net force F acting on the pendulum is the resultant of the
following forces:

1. mg acting downwards.

2. Tension T in the string, acting along its length.

We consider the radial (in the direction of the length of string) and tangential
(in the direction of motion of bob) components of the net force F acting on the
pendulum.
The radial component of force is

Fr = mg cos θ − T (1.92)

The tangential component of force is

Fθ = −mg sin θ (1.93)

23
 1.6 Coupled oscillations

Figure 1.6: A simple pendulum (Source: Physics StackExchange).

Now, the displacement x of the pendulum is always along the tangential direc-
tion, as shown in Figure 1.6. Thus, this tangential component of the force acts
along the direction of displacement and opposes it, as indicated in 1.6. So, this
represents the restoring force acting on the pendulum to cause the oscillatory
motion. Therefore,

d2 x
Fθ = m 2 (1.94)
dt
So, from (1.93) and (1.94),

d2 x
m 2 = −mg sin θ (1.95)
dt
(1.95) does not represent simple harmonic motion (SHM). Recall from (1.4)
that for SHM, the restoring force has to be proportional to the displacement,
but in this case, it is not.
However, if we consider small oscillations, then for small θ,

sin θ ' θ (1.96)

24
 1.6 Coupled oscillations

Therefore, only for small θ,

Fθ = −mgθ (1.97)
x
Also, for small θ, θ ' tan θ ' . Therefore,
L
x
θ' (1.98)
L
Therefore, using (1.96) and (1.98) in (1.95), we get

d2 x x
m 2 = −mg (1.99)
dt L

Rearranging, we get

d2 x g
+ x=0 (1.100)
dt2 L
which is the differential equation representing simple harmonic motion of a
simple pendulum for small oscillations.

Comparing (1.100) with (1.4), we note that

r
2 g g
ω = or ω = (1.101)
L L
and hence, the corresponding time period of oscillations is
s
2π L
T = = 2π (1.102)
ω g

1.6.2 Coupled oscillations - coupled pendulums

Consider two identical pendulums of length l and supporting mass m, coupled
by a weightless spring of spring constant k and natural length equal to the
separation of the masses at zero displacement, as shown in Figure 1.7.

1. Equation of motion
Considering only small oscillations, equations of motion of the two masses,

25
 1.6 Coupled oscillations

Figure 1.7: Two identical pendulums coupled via a massless spring (Source:
Matt Jarvis, Professor of Astrophysics and Fellow, St Cross Col-
lege).

with displacements x and y, respectively, are

d2 x x
m 2 = −mg − k(x − y) (1.103)
dt L

d2 y y
and, m 2 = −mg − k(y − x) (1.104)
dt L
or,

d2 x 2 k
+ ω 0 x = − (x − y) (1.105)
dt2 m
and
d2 y 2 k
+ ω 0 y = − (y − x) (1.106)
dt2 m

where the additional restoring force terms −k(x − y) and −k(y

r− x) are
g
due to coupling of the pendulums via the spring, and ω0 = is the
l
natural frequency of each pendulum.

26
 1.6 Coupled oscillations

Adding (1.105) and (1.106), we get

d2
2
(x + y) + ω02 (x + y) = 0 (1.107)
dt

Subtracting (1.106) and (1.105), we get

d2 2 2k
(x − y) + ω 0 (x − y) = − (x − y)
dt2 m
d2
 
2 2k
or, (x − y) + ω 0 + (x − y) = 0 (1.108)
dt2 m

Replacing

X = x + y in (1.107)
and Y = x − y in (1.108) (1.109)

we get

d2 X
2
+ ω02 X = 0 (1.110)
dt
and
d2 Y
 
2k
2
+ ω02 + Y =0
dt m
d2 Y
or, 2
+ ω12 Y = 0 (1.111)
dt
(1.110) and
s(1.111) represent SHM with natural angular frequencies ω0

2k
and ω1 = ω02 + , respectively.
m

2. Solutions and resultant displacement

General solutions of (1.110) and (1.111), respectively, can be written as1

X = x + y = X0 cos(ω0 t + φ0 ) (1.112)

and

Y = x − y = Y0 cos(ω1 t + φ1 ) (1.113)

1
Note that we can choose to write the general solutions in terms of either sine or cosine functions. Here we
have opted for cosine functions, but the analysis holds equally good with sine functions.

27
 1.6 Coupled oscillations

Based on the relation between x and y, different modes of coupled oscil-

lation may be observed, the most important of which are listed below.

1. In-phase oscillations
In this case x = y, or Y = x − y = 0. Both pendulums follow (1.110)
and oscillate with the angular frequency ω0 . The pendulums are always in
phase. The spring is always unstretched / uncompressed and maintains its
natural length. The relative motion of the pendulums and corresponding
time variation of x and y are schematically represented in Figure 1.8.

Figure 1.8: In-phase coupled oscillations (Sources: Matt Jarvis, Professor of As-
trophysics and Fellow, St Cross College and PPLATO@University
of Reading).

2. Out-of-phase oscillations
In this case x = −y, or X = x + y = 0. Both s pendulums follow (1.111)
 
2k
and oscillate with the angular frequency ω1 = ω02 + > ω0 . The
m
pendulums are always out of phase. The spring is always stretched /

28
 1.6 Coupled oscillations

compressed. The relative motion of the pendulums and corresponding

time variation of x and y are schematically represented in Figure 1.9.

Figure 1.9: Out-of-phase coupled oscillations (Sources: Matt Jarvis, Pro-

fessor of Astrophysics and Fellow, St Cross College and
PPLATO@University of Reading).

These two modes are called normal modes of oscillation. ω0 and ω1 are
called normal frequencies. Any other frequency is a linear combination
of ω0 and ω1 .

3. Resonance
In general, from (1.112) and (1.113),
1
x = (X + Y )
2
1
and, y = (X − Y ) (1.114)
2

29
 1.6 Coupled oscillations

Consider the special condition where the resultant amplitudes and phases
given by (1.112) and (1.113) are equal. That is,

X0 = Y0 = a (say)
and φ0 = φ1 = 0 (the equal value can be zero without loss of generality)
(1.115)

Then, the displacement of the first pendulum

1
x = (X + Y )
2
= a cos ω0 t + a cos ω1 t
(ω1 − ω0 )t (ω1 + ω0 )t
= 2a cos cos (1.116)
2 2
and the displacement of the second pendulum
1
y = (X − Y )
2
= a cos ω0 t − a cos ω1 t
(ω1 − ω0 )t (ω1 + ω0 )t
= 2a sin sin (1.117)
2 2
Clearly, the resultant displacements x and y are a superposition of two
different frequencies; a higher frequency
(ω1 + ω0 )
ωC = (1.118)
2
with a corresponding shorter time period
4π
TC = (1.119)
ω 1 + ω0
and a lower frequency
(ω1 − ω0 )
ωB = (1.120)
2
with a corresponding longer time period
4π
TB = (1.121)
ω1 − ω0

ωB is called Beats frequency.

30
 1.6 Coupled oscillations

Figure 1.10: Resonance in coupled oscillations (Source: PPLATO@University

of Reading).

We can also define the In-phase time period

2π
T0 = (1.122)
ω0
and the Out-of-phase time period
2π
T1 = < T0 (1.123)
ω1

Then from (1.119), (1.121), (1.122), and (1.123), we have

T0 T1
TC = 2 (1.124)
T0 + T1
and
T0 T1
TB = 2 (1.125)
T0 − T1

Degree of coupling is defined as

ω12 − ω02 T02 − T12

χ= 2 = 2 (1.126)
ω1 + ω02 T0 + T12

With the initial conditions x = 2a, y = 0 at t = 0, the displacements of

31
 1.6 Coupled oscillations

the coupled pendulums are as schematically represented in Figure 1.10.

In case of resonance, there is continuous transfer of energy between the
two coupled oscillators, as indicated by the complementary changes in
their respective displacements in Figure 1.10.

Figure 1.11: Spring-mass coupled oscillator system (Source: OSU).

Exercise: Work out the details of coupled oscillations for a spring-mass

system, as shown in Figure 1.11, instead of coupled pendulums discussed above.
That is, set up the equations of motion of the spring-coupled masses, and
discuss the normal modes of oscillations. And thus verify that similar coupled
oscillation features are observed in both systems.

1.6.3 Oscillations in practice

The simple harmonic oscillator is only a mathematical construct to enable un-
derstanding of the physics of oscillations. Practically all oscillators in real-life
are damped or forced.
The most common example of an underdamped oscillator is the simple pen-
dulum. Another good example is a mass-spring system. The swing (as seen in
parks) involves forced oscillations (by someone pushing it). Once the external
force is removed, its motion is damped (just like the pendulum).
Overdamping is typically used in door dampers (and similar dampers) as the
system goes to equilibrium (without overshooting it) over a sufficiently long
time. A shock absorber is essentially a damped spring oscillator, the damping
is from a piston moving in a cylinder filled with oil. If the oil is really thick, or
the piston too tight, the shock absorber will be too stiff - it won’t absorb the
shock, and you will! This is the case of overdamping.
We need to tune the damping so that the car responds smoothly to a bump
in the road, but does not continue to bounce after the bump. this is achieved
using critical damping.
Also note that although we have confined our discussion to mechanical os-
cillations only, oscillations are not only of mechanical nature, but can also be

32
 1.7 Waves

electromagnetic. Mechanical oscillations are characterized by alternating con-

versions of the kinetic energy into one (or several) kinds of potential energy and
back. In electromagnetic oscillations, alternating conversions occur between the
electric field energy (which is analogous to the potential energy in mechanical
systems) and the magnetic field energy (the analogue of the kinetic energy).
Sometimes oscillations have a combined mechanical and electromagnetic na-
ture, e.g., oscillations in plasma.

1.7 Waves
1.7.1 Introduction
Wave motion is periodic motion in space and time. A disturbance of any
conventional property y of a medium may spread through space with time.
Such a spreading of disturbances in a medium is called a wave. y may be
called the wave function.
Let y represent displacement of particles in the medium. Then wave motion
can result from periodic displacement of particles from their mean positions.
The state of motion of a particle is called phase. Different particles may be
in different phases at a given time.

Wave motion consists of the propagation of phase from point to point in the
medium, distinct from the motion of particles in the medium2 .

Surfaces of constant phase in a medium are called wavefronts. A wave

whose wavefronts are infinite parallel planes are called plane waves.

Types of waves:

1. Transverse waves are waves in which the displacement of particles about

their mean position is perpendicular to the direction of propagation. Ex-
amples of transverse waves are EM (electromagnetic) waves and waves in
a pulled string.

2. Longitudinal waves are those in which the displacement of particles

about their mean position is along the direction of propagation. Examples
of longitudinal waves are sound waves and propagation of compressions
and rarefactions in a slinky.

33
 1.7 Waves

Figure 1.12: (a) Longitudinal and (b) transverse waves demonstrated using a
slinky.

In the next Sections, we will focus on transverse waves; discuss the important
parameters, and set up the general wave equation and its differential form.

1.7.2 Parameters associated with a transverse wave

Figure 1.13: Amplitude, wavelength, and time period of a simple harmonic

transverse sinusoidal wave (Source: sengpielaudio).

1. Amplitude (a): maximum displacement of a particle from its mean

position.
2. Time period (T ): time taken to complete one wave cycle, as shown in
Figure 1.13.

34
 1.7 Waves

1
3. Frequency (n or ν): number of wave cycles per second. n = .
T
2π
4. Angular frequency (ω): ω = 2πn = .
T
5. Wavelength (λ): distance travelled by the wave in one cycle, i.e., in
time T , as shown in Figure 1.13.

6. wave number (k): a measure of the number of wave cycles per unit
2π
distance. k = .
λ
λ 1
7. Velocity (v): v = . Since n = , we have
T T
v =n·λ (1.127)

1.7.3 Wave equation (transverse wave)

A simple harmonic progressive wave advances continuously in a given direction
without change in form and particles of the medium perform simple harmonic
oscillations about their mean position with the same amplitude and time period.

Figure 1.14: A simple harmonic progressive wave (Source: AnkPlanet).

We consider a simple harmonic progressive transverse wave as schematically

represented in Figure 1.14. At any time t, the displacements of a particle at
the origin O is given by

y = a sin ωt (1.128)

35
 1.7 Waves

Consider another particle at P , at a distance x from O. Then the displacements

at P are given by

y = a sin(ωt − φ) (1.129)

where φ is the phase difference between O and P .

Now, for one complete wave cycle, i.e., for a particle at B, at a distance λ
from O as shown in Figure 1.14,

x=λ
and φ = 2π (1.130)

because over a complete wave cycle, the phase goes through a cycle of 0 to 2π.
Therefore, from (1.130),
2πx
φ= = kx (1.131)
λ
If v is the wave velocity, then
2πv
ω = 2πn = (1.132)
λ
Using these expressions for φ and ω in (1.129), we can write
 
2πvt 2πx
y = a sin − (1.133)
λ λ
or,

2π
y = a sin (vt − x) (1.134)
λ

For a particle at a distance x in the negative direction from O, then displace-

ment is given by

2π
y = a sin (vt + x) (1.135)
λ

We can rewrite (1.134) and (1.135) in terms of ω and k as

36
 1.7 Waves

y = a sin(ωt − kx) (for positive x) (1.136)

and

y = a sin(ωt + kx) (for negative x) (1.137)

1.7.4 Differential equation of wave motion

The displacement for a simple harmonic progressive transverse wave is given by
(1.136)

y = a sin(ωt − kx) (1.138)

Differentiating3 (1.138) w.r.t. t, twice, we get

∂y
= aω cos(ωt − kx)
∂t
∂ 2y
=⇒ 2
= −aω 2 sin(ωt − kx) (1.139)
∂t
Differentiating (1.138) w.r.t. x, twice, we get
∂y
= −ak cos(ωt − kx)
∂x
∂ 2y
=⇒ 2
= −ak 2 sin(ωt − kx) (1.140)
∂x
Comparing (1.139) and (1.140), we have

∂ 2y ω2 ∂ 2y
= 2· 2 (1.141)
∂t2 k ∂x
Now,
ω 2πn
= =n·λ=v (1.142)
k 2π/λ

Therefore, from (1.141) and (1.142),

3
Note that here we switch to partial derivatives from total derivatives. This is because, unlike in case of
oscillations, where the displacement is a function of time (t) only, in case of waves the displacement is a
function of both time (t) and position (x).

37
 1.7 Waves

∂ 2y 2
2 ∂ y
=v · 2 (1.143)
∂t2 ∂x
which is the differential equation of motion for a transverse wave travelling with
velocity v.

38
 1.7 Waves

Note: Initial phase at the origin and other equivalent forms

Note that, in the above treatment, we have considered that the phase of the
particle at the origin is zero. This does not affect the general applicability of
the wave equation (1.134) or (1.136). However, one may consider some non-
zero initial phase φ0 at the origin, and in this case, the further generalized wave
equation is

y = a sin(ωt − kx − φ0 ) (1.144)

Further note that, instead of (1.134) and (1.136), we can also write the wave
equation as
2π
y = a sin (x − vt) = a sin(kx − ωt) (1.145)
λ
.

1.7.5 Numerical problems

Example 1: The equation of motion of a wave travelling in a spring is

y = 4 sin π(0.10x − 2.0t). Find the (i) amplitude, (ii) wavelength, (iii) initial
phase at the origin, (iv) speed, and (v) frequency of the wave. Here time is
given in seconds and distance in cm.

Solution:

To determine the required parameters, we must rewrite the given equation

in a form similar to the general equation and compare term-by-term.
The general wave equation similar in form to the given equation is
2π
y = a sin (x − vt)
λ
We rearrange the given equation as

y = 4 sin[(0.1π)(x − 20t)]

Comparing term-by-term:
(i) amplitude, a = 4 cm.
2π
(ii) = 0.1π =⇒ wavelength λ = 20 cm.
λ
(iii) Initial phase φ0 = 0.

39
 1.7 Waves

(iv) v = 20 cm/s.
v 20
(v) Frequency n = = = 1.0 Hz.
λ 20

Example 2: A radio station broadcasts at a frequency of 15 × 106 Hz.

Calculate the length of the radio waves, given that velocity of radio waves is
equal to the velocity of light, i.e., 3 × 108 m/s.

Solution:

Given, frequency of the radio waves, n = 15×106 Hz, and velocity v = 3×108
m/s.
Therefore, length of the waves (i.e., the wavelength)
v 3 × 108
λ= = = 20 m.
n 15 × 106
********************* The End (of this Unit) *********************

40
 Interference of light
PH-1007 (Physics)
Dr. Gorky Shaw

2.1 Introduction
Light is a transverse, electromagnetic wave that is visible to the human eye.
Light is dualistic in nature.
Certain optical phenomena, such as interference, diffraction, and polarization,
are explained by the wave nature of light. Light-matter interaction in processes
of emission and absorption, such as photoelectric effect and Compton effect,
are explained by the corpuscular (particle) nature of light.

2.1.1 A few terms

1. Huygens’ principle
Every point on a wavefront is a source of wavelets, which spread forward
at the same speed.

Figure 2.1: Schematic representation of Huygens’ principle (a) Secondary

wavelets in a plane wavefront (b) Secondary wavelets emanating
from a small opening, acting as a point source (Sources: OSU and
HyperPhysics ).

2. Superposition principle
When two waves interact, the resulting wave function is the sum of the
two individual wave functions.
 2.1 Introduction
Consider two waves with wave functions y1 and y2 , represented by the
equations

y1 = a1 sin(ω1 t − k1 x)
and y2 = a2 sin(ω2 t − k2 x) (2.1)

When these two waves interact, the resultant wave function y is given by

y = y1 + y2
= a1 sin(ω1 t − k1 x) + a2 sin(ω2 t − k2 x) (2.2)

In general, for n number of interacting waves, the resultant wave function

is given by

y = y1 + y2 + y3 + ... + yn
Xn
= yi (2.3)
i=1

3. Coherent sources
Two light sources are said to be coherent if they produce waves that have
the same frequency, and a sharply defined phase difference that remains
unchanged with time. The most common example of coherent sources is
Laser.

4. Interference
When two waves produced by coherent sources travel simultaneously in
a medium and superpose each other, the resultant intensity is not dis-
tributed uniformly in space. This modification in intensity is called in-
terference.

5. Constructive and destructive interference

At points where a crest (trough) of one wave falls on a crest (trough) of
the other, the resultant amplitude is the sum of the amplitudes of each
wave separately. The intensity at these points is maximum. This is the
case of constructive interference.
At points where a crest of one wave falls on a trough of the other, the
resultant amplitude is the difference of the amplitudes of each wave. The
intensity at these points is minimum. This is the case of destructive
interference.

2
 2.1 Introduction

Figure 2.2: (a) Constructive and (b) destructive interference, a direct conse-
quence of the superposition principle (Source: A-Level Physics).

2.1.2 Conditions for the interference of light

For sustained interference:
1. The sources must be coherent, i.e., they must maintain a constant phase
difference between them.
2. The two waves must have the same frequency (and hence the same wave-
length). This also means the sources should be monochromatic, i.e., of a
single wavelength.
3. Polarized waves must be in the same state of polarization.
For observation:
1. The separation between the sources should be small.

3
 2.2 Analytical treatment of interference by division of wavefront - Young’s
double slit
2. The distance of observation should be large.
For good contrast:
1. Amplitudes of the interfering waves should be nearly equal (and ideally,
exactly equal).

2. The sources must be narrow.

2.1.3 Types of interference

1. Division of wavefront by employing mirrors, biprisms, or lenses. Divi-
sion of wavefront requires a point source or narrow slit source.
Examples: Fresnel biprism, Young’s double slits, Fresnel’s double mirror,
Lloyd’s mirror.

2. Division of amplitude of incident beam into two or more parts, either

by partial reflection or refraction, which traverse different paths and are
finally brought together to produce interference.
Examples: Interference in thin films, Newton’s rings, Michelson’s inter-
ferometer.

2.2 Analytical treatment of interference by division of

wavefront - Young’s double slit
The wave nature of light was strongly inferred by an experiment performed by
Thomas Young. It was based on how light waves passing through two closely
spaced slits interact. Figure 2.3 shows a schematic representing the processes
in his experiment illustrating interference by division of wavefront. Next we
discuss the observations analytically.
As shown in Figure 2.4, let S be narrow slit illuminated by a monochromatic
light source, and S1 and S2 two similar parallel slits close to each other and
equidistant from S. Then S1 and S2 become centres of secondary wavelets and
act as coherent sources.

2.2.1 Resultant amplitude and intensity

Let a1 and a2 be the amplitudes at P due to the waves from S1 and S2 , respec-
tively. The waves reaching P will have different path lengths S1 P and S2 P.
Hence, path difference between the two waves reaching P from S1 and S2

= S2 P − S1 P (2.4)

4
 2.2 Analytical treatment of interference by division of wavefront - Young’s
double slit

Figure 2.3: Young’s double slit experiment (Source: University of Texas).

Corresponding phase difference at P,

2π
δ= × (S2 P − S1 P) (2.5)
λ
where λ is the wavelength of light used.
We may write the individual displacements at P, with a phase difference of δ
between them as

y1 = a1 sin ωt
and y2 = a2 sin(ωt + δ) (2.6)

where ω is the common angular frequency of the waves.

By the principle of superposition, the resultant displacement at P,

y = y1 + y2
= a1 sin ωt + a2 sin(ωt + δ)
= a1 sin ωt + a2 (sin ωt cos δ + cos ωt sin δ)
= (a1 + a2 cos δ) sin ωt + (a2 sin δ) cos ωt (2.7)

5
 2.2 Analytical treatment of interference by division of wavefront - Young’s
double slit

Figure 2.4: Schematic for analytical treatment of interference by division of

wavefront.

We make the replacements with new variables R and θ given by

a1 + a2 cos δ = R cos θ
and a2 sin δ = R sin θ (2.8)

Then, from (2.7) and (2.8),

y = R cos θ sin ωt + R sin θ cos ωt

or, y = R sin(ωt + θ) (2.9)

(2.9) indicates two notable features of the resultant wave: its frequency re-
mains the same (angular frequency ω same as that of the source waves), and its
nature remains unaltered (similar forms of the wave equations (2.6) and (2.9)).
Here, R is the resultant amplitude. Hence the resultant intensity is I = R2 .
Squaring and adding the two equations in (2.8), we get

I = R2 = a21 + a22 + 2a1 a2 cos δ (2.10)

6
 2.2 Analytical treatment of interference by division of wavefront - Young’s
double slit

2.2.2 Conditions for maxima and minima

From (2.10), intensity at P is maximum if

cos δ = +1
=⇒ δ = 2nπ, n = 0, ±1, ±2, ... (2.11)

And corresponding path difference

λ
= S2 P − S1 P = × δ = nλ (2.12)
2π
Maximum value of intensity,

Imax = a21 + a22 + 2a1 a2

= (a1 + a2 )2
> a21 + a22 = I1 + I2 (2.13)

That is,

Imax > I1 + I2 (2.14)

where I1 and I2 are the intensities of the source waves at P.

Special case: equal source amplitudes

If a1 = a2 = a and I = a2 , then

Imax = 4a2 = 4I (2.15)

On the other hand, intensity at P is minimum if

cos δ = −1
=⇒ δ = (2n + 1)π, n = 0, ±1, ±2, ... (2.16)

And corresponding path difference

λ λ
= S2 P − S1 P = × δ = (2n + 1) (2.17)
2π 2

7
 2.2 Analytical treatment of interference by division of wavefront - Young’s
double slit

Minimum value of intensity,

Imin = a21 + a22 − 2a1 a2

= (a1 − a2 )2
< a21 + a22 = I1 + I2 (2.18)

That is,

Imin < I1 + I2 (2.19)

where I1 and I2 are the intensities of the source waves at P.

Special case: equal source amplitudes

If a1 = a2 = a and I = a2 , then

Imin = 0 (2.20)

2.2.3 Interference and conservation of energy

Figure 2.5: Variation of resultant intensity (I) with phase difference (δ).

To determine the average intensity of the resultant wave, we need to deter-

mine the average intensity over a single complete cycle, i.e., between δ = 0 and

8
 2.2 Analytical treatment of interference by division of wavefront - Young’s
double slit

δ = 2π, which is given by

R 2π
Idδ
Iavg = R0 2π
0 dδ
R 2π 2 2
0 (a1 + a2 + 2a1 a2 cos δ)dδ
= R 2π
0 dδ
2π(a21 + a22 )
=
2π
= a1 + a22
2

= I1 + I2 (2.21)

Hence, although there is a variation of resultant intensity between Imax and

Imin , the average resultant intensity Iavg is equal to the sum of the intensities
of the source waves. This means, there is no loss or gain of energy due to
interference. There is just a redistribution of the energy from the source waves.

2.2.4 The (not-so-obvious) small angle approximation

In the analytical treatment above, there is no explicit mention of (i) the small
separation between the sources, or (ii) large distance of observation, both of
which are mentioned earlier as necessary conditions for observation of interfer-
ence pattern. These are indeed incorporated, but in a subtle manner.
Expression (2.8) for the resultant displacement at P, based on the superpo-
sition principle, is valid only if the paths S1 P and S2 P traversed by the source
waves are parallel to each other. As evident from Figure 2.4, in general, the two
paths would not be parallel. Only if the separation between the sources is much
smaller than the distance of observation, i.e., if D >> 2d, can we consider the
two paths S1 P and S2 P to be approximately parallel.

9
 2.2 Analytical treatment of interference by division of wavefront - Young’s
double slit

2.2.5 Fringe width in double slit interference

The bright and dark bands produced as a result of interference of light are
known as interference fringes. We may estimate the width of these fringes
as discussed below.
In Figure 2.4, we have

(S2 P)2 = (S2 M2 )2 + (PM2 )2

= D2 + (x + d)2
(x + d)2
 
2
=D 1+
D2
"  2 #1/2
x+d
=⇒ S2 P = D 1 + (2.22)
D

We consider only small distances x << D. And as discussed above, for

observation of interference pattern, we must have d << D. Hence, we have
(x + d) << D. So, in the above expression, we can use binomial expansion and
keep the first two terms of the expansion, ignoring higher order terms1 , we get
"  2 #
1 x+d
S2 P = D 1 +
2 D
(x + d)2
=D+ (2.23)
2D
Similarly,

(x − d)2
S1 P = D + (2.24)
2D
Therefore,
2xd
S2 P − S1 P = (2.25)
D
For maxima or bright fringes:

   
1 n
P∞ n k n n(n − 1)...(n − k + 1)
For a real number n, the binomial expansion of (1+x) = k=0 k x , where k = .
k!
Here k is a positive integer.

10
 2.2 Analytical treatment of interference by division of wavefront - Young’s
double slit

2xd
= nλ
D
Dλ
=⇒ x = n · (2.26)
2d
For minima or dark fringes:

2xd λ
= (2n + 1)
D 2
(2n + 1) Dλ
=⇒ x = · (2.27)
2 2d
Therefore, spacing between two consecutive bright (or dark) fringes

= xn+1 − xn
Dλ
= (2.28)
2d
which is independent of n.

It is called fringe width (β).

Dλ
β= (2.29)
2d

One may define the angular fringe width βθ as

β λ
βθ = = (2.30)
D 2d

11
 2.3 Numerical problems - Interference by division of wavefront

2.3 Numerical problems - Interference by division of

wavefront

Exercise 2.1
Two coherent sources with intensity ratio 100:1 produce interference fringes.
Find the ratio between maximum and minimum intensities in the interference
pattern.

Solution:

We know,

Imax (a1 + a2 )2
=
Imin (a1 − a2 )2
Now, given,

I1 a21 100
= 2=
I2 a2 1
a1
=⇒ = 10
a2
=⇒ a1 = 10a2

Therefore,

Imax (11a2 )2 121

= =
Imin (9a2 )2 81

Exercise 2.2
The ratio between maximum and minimum intensities in a double slit inter-
ference pattern is 36:1. Find the ratio between amplitudes and intensities of
the two interfering waves.

Solution:

12
 2.3 Numerical problems - Interference by division of wavefront

Given,

Imax (a1 + a2 )2 36
= =
Imin (a1 − a2 )2 1
=⇒ a1 + a2 = 6(a1 − a2 )
=⇒ (a1 + a2 ) = 6(a1 − a2 )
=⇒ 5a1 = 7a2

Therefore, ratio of amplitudes of the interfering waves,

a1 7
=
a2 5
and the ratio of their intensities,

I1 a21 49
= 2=
I2 a2 25

Exercise 2.3
In a Young’s double slit experiment with monochromatic light of wavelength
6000 Å, the fringe width is found to be 0.5 mm in the interference pattern on a
screen at a distance 1 m from the slits. Find the separation between the slits.

Solution:

Given, β = 0.5 mm = 0.05 cm,

λ = 6000 Å= 6 × 10−7 cm,
D = 1 m = 100 cm.
Therefore, separation between the slits,

Dλ 100 × 6 × 10−7
2d = = = 0.0012 cm
β 0.05

Exercise 2.4
Two coherent sources placed 0.2 mm apart produce an interference pattern
observed on a screen 1 m away. With a certain monochromatic light source, the
fourth bright fringe is situated at a distance 10.0 mm from the central fringe.
Find the wavelength of light used.

13
 2.3 Numerical problems - Interference by division of wavefront

Solution:

Position of the nth bright fringe,

Dλ
xn = n ·
2d
Here,
n = 4,
D = 1 m = 1000 mm,
2d = 0.2 mm,
xn = 10 mm.
Therefore,
2dxn
λ=
nD
0.2 × 10
=
4 × 1000
= 5 × 10−4 mm
= 5000 Å

Exercise 2.5
In an interference pattern, 12th order maximum is observed for λ = 6000 Å.
What order maximum is visible with light of wavelength 4800 Å?

Solution:

We know, the path difference between source waves for the nth order maxi-
mum ∝ nλ.
Given, maxima (of different orders) for two wavelengths occur at the same
position, or

n2 λ2 = n1 λ1
n1 λ1 12 × 6000
or, n2 = = = 15
λ2 4800

14
 2.4 Interference due to thin films

2.4 Interference due to thin films

2.4.1 Phase change on reflection - Stokes’ treatment
When a light wave is reflected at the surface of an optically denser medium
(with higher refractive index), it suffers a phase change of π (but not when
reflected at the surface of a rarer medium (with lower refractive index)).
This is illustrated in Figure 2.6.

Figure 2.6: Phase change on reflection (Source: Stack Exchange).

2.5 Interference due to parallel-sided thin films: reflected

rays
Let a monochromatic light beam SA be incident at angle i on a parallel-sided
transparent thin film of thickness t and refractive index µ > 1, as shown in
Figure 2.7. At A, it is partly reflected along AR1 at angle i, and partly refracted
along AB at angle r. At B, it is again partly reflected along BC and partly
refracted along BT1 . Similar reflections and refractions occur at C, D etc.
Path difference between the rays AR1 and CR2 ,

p = µ(AB + BC) − AN (2.31)

15
 2.5 Interference due to parallel-sided thin films: reflected rays

Figure 2.7: Interference by division of amplitude due to a parallel-sided thin

film.

But,

AB = BC = BM sec r = t sec r (2.32)

and

AN = AC sin i = (AM + MC) sin i

= 2AM sin i (since AM = MC)
= 2BM tan r sin i (since AM = BM tan r)
= 2t tan r sin i
sin i
= 2t tan r(µ sin r) (since = µ)
sin r
= 2µt sec r sin2 r (2.33)

16
 2.5 Interference due to parallel-sided thin films: reflected rays

Therefore,

p = 2µt sec r − 2µt sec r sin2 r

= 2µt sec r cos2 r
= 2µt cos r (2.34)

Further, the ray AR1 undergoes an additional phase difference of π, and

λ
hence, and additional path difference of , where λ is the wavelength of light
2
used.
Therefore, effective path difference between AR1 and CR2
λ λ
peff = p ± = 2µt cos r ± (2.35)
2 2
Thus, depending on the values of t and λ, the effective path difference between
the reflected rays can lead to constructive or destructive interference, and the
film may appear dark or bright in reflected light.
We can use either of the + and − signs. In the following discussion, we use
the + sign.

2.5.1 Conditions for maxima and minima

For maximum intensity of reflected light,

λ
2µt cos r + = nλ, n = 1, 2, 3...
2
λ
or, 2µt cos r = (2n − 1) (2.36)
2
In this condition, the thin film will appear bright in reflected light.
For minimum intensity of reflected light,
λ λ
2µt cos r + = (2n + 1) , n = 0, 1, 2, 3...
2 2
or, 2µt cos r = nλ (2.37)

In this condition, the thin film will appear dark in reflected light.

Note that the choice of + and − signs, as well as values of n (0, 1, 2... or 1, 2, 3...)
in the conditions (2.36) and (2.37) are not arbitrary. These are carefully chosen

17
 2.6 Interference due to thin films - further points

so that the value of t, the thickness of the film, is always positive. Negative
value of t would not make physical sense.

2.6 Interference due to thin films - further points

2.6.1 Interference due transmitted rays
There is no additional phase change due to reflection at points like B and C in
Figure 2.7 because, in these cases, light is travelling from the denser medium
to the rarer medium.
For maximum intensity of transmitted light,

2µt cos r = nλ (2.38)

For minimum intensity of reflected light,

λ
2µt cos r = (2n − 1) (2.39)
2

2.6.2 Colours in thin films

White light is composed of light of different wavelengths. Since every colour
of light has a different wavelength, according to (2.36), only one colour can
constructively interfere at a given thickness of a film. Hence, different colours
appear brighter for different thicknesses of films.
Some common examples of this phenomenon are the rainbow colour patterns
seen in soap bubbles, oil spills, bottoms of metallic cooking pots etc. (see Figure
2.8).

2.6.3 Interference due to wedge-shaped thin films

Figure 2.9: Wedge-shaped thin film.

18
 2.6 Interference due to thin films - further points

Figure 2.8: Colours in thin films (a) Soap bubbles, (b) Oil spill, and (c) Bottom
of a cooking pot.

For a wedge-shaped thin film with wedge angle θ, as shown in Figure 2.9, the
effective path difference at thickness t is given by
λ
peff = 2µt cos (r + θ) ± (2.40)
2
Accordingly, in reflected light, for maximum intensity,
λ
2µt cos (r + θ) = (2n − 1) , n = 1, 2, 3... (2.41)
2
and for minimum intensity,

2µt cos (r + θ) = nλ, n = 0, 1, 2, 3... (2.42)

For such a film, the thickness varies across the film. Hence, alternate locations
will demonstrate constructive and destructive interference, according to the
thickness of the film at that location. This is the basis of Newton’s rings
formation, which is discussed next.

19
 2.7 Newton’s rings

2.7 Newton’s rings

Newton’s rings is a phenomenon in which interference pattern is created by the
reflection of light from an air film formed between a spherical surface and an
adjacent touching flat surface.
With monochromatic light, a series of concentric alternate bright and dark
rings are observed. With white light, concentric ring patterns of rainbow colours
is observed, because the different wavelengths of light interfere constructively
at different thicknesses of the air film.

Figure 2.10: Experimental setup for observation of Newton’s rings.

A typical experimental setup for observation of Newton’s rings is schemati-

cally represented in Figure 2.10. Usually, monochromatic light is made incident
normally on a combination of a plano-convex lens and a plane glass plate, as
shown in the Figure. The reflected light is observed using a travelling micro-
scope. A thin air film of variable thickness is formed between the lens and the
glass plate.
Alternate bright and dark circular rings, with a central dark spot, are ob-
served in the reflected light. This Newton’s rings formation is a result of in-

20
 2.7 Newton’s rings

terference between light waves reflected from the upper and lower surfaces of
the wedge-shaped air film. The thickness of the air film is zero at the point of
contact and gradually increases outwards. Due to the shape of the plano-convex
lens, the locus of points where the thickness of the air film is constant is a circle,
with the point of contact as its centre.

2.7.1 Effective path difference and the central dark spot

The air film can be considered a wedge-shaped thin film. Then the effective
path difference between waves reflected from the upper and lowers surfaces at
thickness t is given by
λ
peff = 2µt cos (r + θ) + (2.43)
2
Now, for large radius of curvature of the plano-convex lens, θ can be ignored.
Further, for normal incidence, r = 0. Therefore, we have
λ
peff = 2µt + (2.44)
2
Therefore, similar to the conditions in Section 2.5.1, for maximum intensity
(bright rings),
λ
2µt = (2n − 1) , n = 1, 2, 3... (2.45)
2
and for minimum intensity (dark rings),

2µt = nλ, n = 0, 1, 2, 3... (2.46)

At the point of contact, t = 0. Hence, from (2.44), at this location,

λ
peff = (2.47)
2
for all λ. This is the condition for minimum intensity. Hence, the central spot
is always dark.

2.7.2 Diameter of rings

Consider the plano-convex lens LOL’ placed on a glass plate AB, the point of
contact being O, as shown in Figure 2.11. Let R be the radius of curvature of

21
 2.7 Newton’s rings

Figure 2.11: Determination of diameter of Newton’s rings.

the lens. Let rn be the radius, and Dn = 2rn the diameter, of the Newton’s ring
corresponding to a point P where the film thickness is t.
We use the property of a circle2 :

rn2 = PN2 = ON × NE
= t × (2R − t)
= 2Rt − t2 (2.48)

Considering R >> t, we have

rn2 = 2Rt
r2
=⇒ 2t = n (2.49)
R

2
Verify this yourself!

22
 2.7 Newton’s rings

Therefore, from (2.45) and (2.49), for a bright ring (constructive interference),

rn2 λ
µ = (2n − 1) , n = 1, 2, 3...
R 2
λR
=⇒ rn2 = (2n − 1) (2.50)
2µ

That is, for bright rings,

λR
Dn2 = 2(2n − 1)
µ
p
or, Dn ∝ (2n − 1) (2.51)

From (2.46) and (2.49), for a dark ring (destructive interference),

rn2
µ = nλ, n = 0, 1, 2, 3...
R
nλR
=⇒ rn2 = (2.52)
µ

That is, for dark rings,

4nλR
Dn2 =
µ
√
or, Dn ∝ n (2.53)

For air film, µ = 1. Therefore,

2
Dn[bright] = 2(2n − 1)λR
2
and Dn[dark] = 4nλR (2.54)

2.7.3 Determination of wavelength of monochromatic light using

Newton’s rings
From (2.54), for an air film, the diameter of the mth dark ring is given by
2
Dm = 4mλR (2.55)

23
 2.7 Newton’s rings

and that of the (m + p)th ring is

2
Dm+p = 4(m + p)λR (2.56)

Therefore,
2 2
Dm+p − Dm = 4pλR (2.57)

That is,

2 2
Dm+p − Dm
λ= (2.58)
4pR

In the experimental setup shown in Figure 2.10, the cross-wire of the eye-
piece of the microscope is focused on any dark ring (say 24th, 16th etc.), and
reading of the microscope is noted. Similar readings are noted by focusing the
cross-wire after every few dark rings. Diameters of different rings are calculated,
and a graph is plotted between the square of the diameter (Dn2 ) and the ring
number n, which is a straight line, as shown in Figure . Its slope gives the value
2 2
Dm+p − Dm
of . If R is known, the wavelength λ of the monochromatic light
p
can then be determined from 2.58.

Figure 2.12: (Dn2 ) versus n plot.

24
 2.8 Numerical problems - Interference by division of amplitude

2.7.4 Determination of refractive index of a liquid using Newton’s rings

For air film, from (2.57),
2 2
[Dm+p − Dm ]air = 4pλR (2.59)

If a liquid is put in between the plano-convex lens and the glass plate, then for
the liquid film,

2 2 4pλR
[Dm+p − Dm ]liquid = (2.60)
µ
Therefore, from (2.59) and (2.60),
2 2
[Dm+p − Dm ]air
µ= 2 2]
(2.61)
[Dm+p − Dm liquid

Hence, refractive index (µ) of the liquid can be determined by measuring the
diameters of Newton’s rings in the liquid and in air.

2.8 Numerical problems - Interference by division of

amplitude

Exercise 2.6
A soap film (µ = 1.33) seen by sodium light (λ = 5893 Å) by normal reflection
appears dark. Find the minimum thickness of the film.

Solution:

For the film to appear dark, i.e., for destructive interference, we have

2µt cos r = nλ

Clearly, the minimum possible value of the thickness of the film, t, is obtained
for n = 1.
For normal incidence, r = 0, cos r = 1
Therefore,
nλ 1 × 5893
t= = = 2215 Å
2µ 2 × 1.33

25
 2.8 Numerical problems - Interference by division of amplitude

Exercise 2.7
Newton’s rings are observed in reflected light of wavelength 6250 Å. Diameter
of the 10th dark ring is 0.5 cm. Find (i) the radius of curvature of the lens (ii)
thickness of the air film.

Solution:

For air film, diameter of the nth dark ring, Dn , is given by

Dn2 = 4nλR

Given, n = 10, Dn = 0.5 cm, λ = 6250 Å= 6.25 × 10−5 cm.

(i) Therefore, radius of curvature of the lens

Dn2 (0.5)2
R= =
4nλ 4 × 10 × 6.25 × 10−5
0.25
=
250 × 10−5
= 100 cm
= 1m

(ii) Condition for dark ring:

2t = nλ
nλ
=⇒ t =
2
10 × 6.25 × 10−5
=
2
= 31.25 × 10−5
= 3.125 µm

Exercise 2.8
In a Newton’s rings experiment, the diameter of the 15th ring is found to
be 0.590 cm and that of the 5th ring 0.336 cm. If radius of curvature of the
plano-convex lens is 100 cm, calculate the wavelength of light used.

Solution:

26
 2.8 Numerical problems - Interference by division of amplitude

Here, m + p = 15, m = 5. Therefore, p = 15 − 5 = 10.

And, D15 = 0.590 cm, D5 = 0.336 cm, R = 100 cm.
Therefore,
2 2
Dm+p − Dm 0.5902 − 0.3362
λ= =
4pR 4 × 10 × 100
0.3481 − 0.1129
=
4 × 103
= 0.0588 × 10−3 cm
= 5880 × 10−8 cm
= 5880 Å

Exercise 2.9
In a Newton’s rings formation in a liquid film with light of wavelength 6000
Å, diameter of the sixth bright ring is 3.1 mm and radius of curvature of the
curved surface is 1 m. Calculate the refractive index of the liquid.

Solution:

Here n = 6, D6 = 3.1 mm = 0.31 cm, R = 1 m = 100 cm, λ = 6000

Å= 6 × 10−5 cm.
For bright rings,
λR
Dn2 = 2(2n − 1)
µ
2(2n − 1)λR
=⇒ µ =
Dn2
Therefore, in this case,

2 × (12 − 1) × 6 × 10−5 × 1
µ=
0.31
2 × 11 × 6 × 10−5
=
0.31
≈ 1.374

27
 2.9 Michelson interferometer

2.9 Michelson interferometer

The Michelson interferometer produces interference fringes by splitting a light
beam into two parts and then recombining them after they have travelled dif-
ferent optical paths. It has numerous applications, including precise length
measurement.
As shown in Figure 2.13, light from a diffuse light source S is made incident
on a half-silvered mirror M which passes half of the light and reflects the other
half to the mirror M1 . The passed (refracted) light is then made to pass through
a fully transparent compensation plate C onto a second mirror M2 . The com-
pensation plate is made from same material and has the same thickness as M.
This ensures both beams pass through the same thickness of glass, and any
phase difference between the two beams is only due to the difference in the
distances they travel.

Figure 2.13: Michelson interferometer (a) Typical experimental setup (b) Pla-
nar view showing the propagation of light in the set up (Source:
Physics LibreTexts).

The observer sees the interference pattern due to beams reflected from the
mirror M1 and a virtual image M02 of the mirror M2 formed by M, which act
as coherent sources. The arrangement is optically equivalent to a thin air film
between M1 and M02 .
The schematics in Figure 2.13 only a single reference light beam emergent
from the light source S. However, it is an extended source (rather than a point
source), which means there are light beams are emergent from various points
on the source and are incident on M at different angles. With monochromatic
light, the resultant of interference between all these beams from the extended
source is a circular pattern of alternating dark and bright fringes (similar in
appearance to Newton’s rings) when the two mirrors M1 and M02 are parallel

28
 2.9 Michelson interferometer

to each other. Each circular ring corresponds to a certain inclination of the

incident beam with respect to the plane of the mirrors. Hence the fringes are
known as fringes of equal inclination. When M1 and M02 are not parallel, they
enclose a wedge-shaped film, and curved fringes can be observed. These are
also known as fringes of equal thickness. When M1 and M02 intersect, straight
line fringes are obtained around the point of intersection.

2.9.1 Suggested further reading: Michelson interferometer

1. Physics LibreTexts

2. Physical Optics Blog

3. Michelson interferometer@Physics StackExchange

4. Carl Wieman Science Education Initiative (CWSEI)

5. Michelson Interferometer Experiment manual@NISER

...and a lot more on the World Wide Web.

********************* The End (of this Unit) *********************

29
 Diffraction
PH-1007 (Physics)
Dr. Gorky Shaw

3.1 Introduction
According to geometrical optics, if a plane wave is incident on a long narrow
slit, as shown in Figure 3.1, the region PQ is illuminated and the rest of the
region forms the geometrical shadow.
However, if the slit is made very narrow (comparable to the wavelength of
light), then light bends into the geometrical shadow region. This is known as
diffraction.

Figure 3.1: Light bent into geometrical shadow by diffraction.

1. Diffraction is the slight bending of light into the geometrical shadow

region, when passing through a narrow slit, comparable to the wavelength
of light.
2. The resultant intensity distribution is known as diffraction pattern.
3. Diffraction can be explained based on Huygens principle which states that
each point on an wavefront acts as a source of secondary waves, called
wavelets. These superpose with each other to result in the phenomenon
of diffraction.
4. Diffraction results from the interference between waves from the same
source.
 3.2 Fraunhofer diffraction at a single slit
Some differences between interference and diffraction are:

1. Interference results from waves emerging from different sources. Diffrac-

tion results from waves emerging from the same source.

2. Diffraction maxima are never of the same width. Interference maxima

may or may not be of the same width.

3. Diffraction requires an obstacle in the path of light. Interference can occur

without any obstacle.

4. Direction propagation remains unaltered in interference, whereas light

gets bent and direction of propagation is modified in diffraction.

3.1.1 Types of diffraction

Diffraction can be categorized into two types:

1. Fresnel diffraction: The source, slit and screen are at finite distance
from each other. The incident wavefront is spherical or cylindrical.

2. Fraunhofer diffraction: The source and screen are effectively at infinite

distance from the diffracting obstacle. This is achieved by placing the
source and the screen at the focal planes of two lenses. The incident
wavefront is plane.

3.2 Fraunhofer diffraction at a single slit

3.2.1 Resultant of n simple harmonic waves
For n simple harmonic waves having equal amplitude a and common phase
difference d between successive vibrations, the resultant amplitude R is given
by
a sin(nd/2)
R= (3.1)
sin(d/2)

3.2.2 Resultant intensity due to Fraunhofer diffraction at a single slit

We consider Fraunhofer diffraction from a single slit AB as shown in Figure
3.2. According to Huygens principle, each point on slit AB sends out sec-
ondary wavelets in all directions. Rays which do not diffract focus at O. Those
diffracted through an angle θ are focused at P. We have, path difference between

2
 3.2 Fraunhofer diffraction at a single slit

Figure 3.2: Fraunhofer diffraction at a single slit.

wavelets reaching P from A and B,

BK = AB sin θ = e sin θ (3.2)

Corresponding phase difference

2π
= e sin θ (3.3)
λ
Let the width of AB be divided into n equal parts, and amplitude of wavelets
from each part be a.
Then, the phase difference between waves from consecutive parts is
 
1 2π
d= e sin θ (3.4)
n λ

Therefore, the resultant amplitude at P,

a sin(nd/2)
R=
sin(d/2)
a sin(πe sin θ/λ)
= (3.5)
sin(πe sin θ/nλ)

3
 3.2 Fraunhofer diffraction at a single slit

πe sin θ
Let = α. Then
λ
a sin(nd/2)
R=
sin(d/2)
a sin α
=
sin(α/n)
a sin α
' (for small α/n)
α/n
na sin α
= (3.6)
α
As n → ∞, a → 0. But the product na remains finite.
Let na = A. Then
A sin α
R= (3.7)
α
Therefore, resultant intensity at P,

A2 sin2 α
2
I=R = (3.8)
α2

3.2.3 Conditions for maxima and minima

For maximum or minimum intensity, we must have
dI
=0
dα
d A2 sin2 α
 
=⇒ =0
dα α2
2 sin2 α
 
2 2 sin α cos α
=⇒ A − =0
α2 α3
=⇒ sin α(α cos α − sin α) = 0 (3.9)

This means, either

sin α = 0 (3.10)

4
 3.2 Fraunhofer diffraction at a single slit

or

α cos α − sin α = 0
or, α − tan α = 0 (3.11)

3.2.4 Condition for minima

From (3.8) and (3.10), the intensity at P is minimum (Imin = 0) if

sin α = 0 (but α 6= 0)
=⇒ α = nπ (n = ±1, ±2, ±3...but n 6= 0)
πe sin θ
=⇒ = nπ (3.12)
λ
That is,

e sin θ = nλ (3.13)

First, second, third... minima occur at n = ±1, ±2, ±3... etc.

3.2.5 Principal maximum

sin α
When α = 0, the term becomes indeterminate. Clearly, this is not a
α
condition for minimum intensity. So, this has to be condition for maximum
intensity. To establish this, we consider the limit α → 0 and apply L’Hospital’s
rule as follows:
sin α cos α
lim = lim = 1 6= 0 (3.14)
α→0 α α→0 1
So, α = 0 does correspond to a maximum. Thus, we have

πe sin θ/λ = 0
=⇒ sin θ = 0
=⇒ θ = 0 (3.15)

That is the same direction as that of the incident light. This is known as the
principal maximum or central maximum.

5
 3.2 Fraunhofer diffraction at a single slit

3.2.6 Secondary maxima

Since there are a number of minima (given by the condition (3.13), there must
also be maxima in between any two of these minima). These are known as
secondary maxima, and the condition for these is given by (3.11), that is

α − tan α = 0 (3.16)

(3.16) can be solved graphically by plotting y = α and y = tan α and finding

the points of intersection. We can thus obtain the solutions of (3.16):
3π 5π 7π
α'± , ± , ± , ... (3.17)
2 2 2
apart, of course, from α = 0, which is the condition for the central maximum.
From (3.17), we have for secondary maxima,

πe sin θ π
' (2n + 1) , n = ±1, ±2, ±3... (3.18)
λ 2
or,

λ
e sin θ ' (2n + 1) , n = ±1, ±2, ±3... (3.19)
2

3.2.7 Ratio of intensities at the principal and secondary maxima

Intensity at the central maximum,
 2
sin 0
I0 = A2 = A2 (3.20)
0

At the secondary maxima,

2
4A2

sin(3π/2)
I1 = A2 = (for n = 1)
3π/2 9π 2
2
4A2

2 sin(5π/2)
I2 = A = (for n = 2)
5π/2 25π 2
2
4A2

2 sin(7π/2)
I3 = A = (for n = 3) (3.21)
7π/2 49π 2

6
 3.2 Fraunhofer diffraction at a single slit

etc.

So, intensities at successive maxima are in the ratio

4A2 4A2 4A2

1: : : ... (3.22)
9π 2 25π 2 49π 2
as represented schematically in Figure 3.3.

Figure 3.3: Intensity distribution due to Fraunhofer diffraction at a single slit.

7
 3.3 Plane Diffraction Grating

3.3 Plane Diffraction Grating

A plane diffraction grating is an arrangement having a large number of
parallel slits of the same width, separated by opaque spaces of equal width.

3.3.1 Resultant intensity

Let AB be the section of a grating with slit width e and the width of opaque
spaces d (as shown in Figure 3.4). The quantity (e + d) is called the Grating
element. Let N be the number of slits in the grating. Path difference between
rays from slits S1 and S2 at P ,

S2 K = S1 S2 sin θ = (e + d) sin θ (3.23)

Corresponding phase difference,

2π
= (e + d) sin θ = 2β (say) (3.24)
λ
where λ is wavelength of the incident monochromatic light.
Then, in direct analogy with the case of diffraction from a single slit, the
resultant amplitude at P is
0
A sin N β
R= (3.25)
sin β

0
where A = Resultant amplitude from each slit
A sin α
= (3.26)
α
That is,

A sin α sin N β
R= · (3.27)
α sin β

So, resultant intensity at P ,

0
A 2 sin2 N β
2 A2 sin2 α sin2 N β
I=R = = · (3.28)
sin2 β α2 sin2 β

8
 3.3 Plane Diffraction Grating

The first factor gives a diffraction pattern due to a single slit. The second
factor gives the interference pattern due to N slits.

Figure 3.4: Schematic representation of a plane diffraction grating.

3.3.2 Maxima and minima in the diffraction pattern

For maximum or minimum value of intensity I given in (3.4), we have the
condition
dI
=0
dβ
 0
d A 2 sin2 N β

=⇒ =0
dβ sin2 β
2N sin N β cos N β 2 sin2 N β cos β
=⇒ − =0
sin2 β sin3 β
=⇒ sin N β(N cos N β sin β − sin N β cos β) = 0

The above implies either

9
 3.3 Plane Diffraction Grating

sin N β = 0 (3.29)

or

N cos N β sin β − sin N β cos β = 0

Dividing by cos N β cos β and rearranging, we get

N tan β = tan N β (3.30)

3.3.3 Minima
When sin N β = 0 (but sin β 6= 0), we have the minimum intensity I = Imin = 0.
sin N β
When sin β = 0, the term is indeterminate (and corresponds to principal
sin β
maxima, as discussed in Section 3.3.4 below.
Therefore, the condition for minimum intensity, given by (3.29) is

sin N β = 0( but sin β 6= 0)

=⇒ N β = mπ (3.31)

π
=⇒ N (e + d) sin θ = mπ
λ

=⇒ N (e + d) sin θ = mλ (3.32)

where m takes all integral values except 0, ±N, ±2N , ±3N , ..., ±nN (where
n is an integer). This is because these values of m give sin β = 0, corresponding
to principal maxima, as discussed in Section 3.3.4. Thus, there are (N − 1)
minima between successive principal maxima.

10
 3.3 Plane Diffraction Grating

3.3.4 Principal maxima

sin N β
For maximum intensity, should be maximum. As discussed above, this
sin β
expression is indeterminate, and the intensity is maximum for for sin β = 0, or,

β = nπ, n = 0, ±1, ±2, ±3... (3.33)

Therefore we take limits and apply L’Hospitals’ rule:

sin N β N cos N β
lim = lim =N (3.34)
β→nπ sin β β→nπ cos β
Corresponding intensity,

A2 sin2 α
Ip = 2
· N2 (3.35)
α

These maxima are most intense and are called principal maxima. These
are obtained in the directions given by the condition (3.33):

β = nπ
π
=⇒ (e + d) sin θ = nπ
λ

=⇒ (e + d) sin θ = nλ (3.36)

where n = 0, ±1, ±2, ±3... correspond to the nth order principal maxima.

3.3.5 Secondary maxima

The (N −1) minima between each principal maximum produce (N −2) maxima,
called secondary maxima. These can be obtained from the condition (3.30),
sin N β
i.e., N tan β = tan N β. To find the corresponding values of , we use the
sin β
triangle as shown in Figure 3.5 which depicts the condition (3.30).

11
 3.3 Plane Diffraction Grating

sin N β
Figure 3.5: Triangle to determine .
sin β

From the triangle, we can see that

N tan β
sin N β = p (3.37)
1 + N 2 tan2 β
Therefore,

sin2 N β N 2 tan2 β
=
sin2 β (1 + N 2 tan2 β) sin2 β
N2
=
(cot2 β + N 2 ) sin2 β
N2
=
cos2 β + N 2 sin2 β
N2
= 2 2 (using sin2 β + cos2 β = 1) (3.38)
1 + (N − 1) sin β
Therefore, intensity of the secondary maxima,

0 A2 sin2 α N2
I = ·
α2 1 + (N 2 − 1) sin2 β

0 Ip
=⇒ I = (3.39)
1 + (N 2 − 1) sin2 β

The intensity of secondary maxima relative to principal maxima decreases

with increasing N .

12
 3.3 Plane Diffraction Grating

sin2 α sin2 N β
Variation of the and terms with the direction θ, and the
α2 sin2 β
resultant grating spectrum (intensity curve), are shown in Figure 3.6.

Figure 3.6: Grating spectrum

3.3.6 Missing orders or absent spectra

Missing orders or absent spectra occur occur when a principal maximum in
the grating spectrum overlaps with a minimum in the corresponding single slit
diffraction pattern.
From (3.36), principal maximum in the diffraction spectrum occur when

(e + d) sin θ = nλ, n = 0, ±1, ±2, ±3... (3.40)

Now, minima in a single slit pattern occur when

e sin θ = mλ, m = ±1, ±2, ±3... (3.41)

If both (3.40) and (3.41) are satisfied for a given θ, a particular maximum of
order n will be missing in the grating spectrum.

13
 3.4 Dispersive power and resolving power of a grating

Comparing (3.40) and (3.41), we have,

e+d n
= (3.42)
e m
Thus, for the nth order spectrum to be absent, for d = e,

n = 2m. i.e., n = 2, 4, 6... for m = 1, 2, 3... (3.43)

3.4 Dispersive power and resolving power of a grating

3.4.1 Dispersive power
The dispersive power of a grating
 is defined as the rate of variation of angle
dθ
of diffraction with wavelength .
dλ
Differentiating the condition for principal maxima for a given wavelength
with respect to λ, we get
dθ
(e + d) cos θ =n
dλ

dθ n
=⇒ = (3.44)
dλ (e + d) cos θ

3.4.2 Resolving power of an optical instrument: Rayleigh’s criterion of

resolution
Two images are just resolved if the position of the principal maximum in the
diffraction pattern of one coincides with the first minimum in the diffraction
pattern of the other. The criterion is illustrated in Figure 3.7. In panel (a), the
two neighbouring principal maxima corresponding to two wavelengths λ1 and
λ2 are clearly distinguishable due to the fairly wide region of low intensity in
between them. In panel (b), the two principal maxima are too close to each
other to be distinguished, as the resultant intensity shows no drop in between
them, and shows a peak instead. It would be difficult to ascertain whether the
pattern corresponds to two lines or just a single line. In panel (c), the principal
maximum of one overlaps exactly with the first minimum of the other, and
there is a dip in intensity in between the two, so that they are just resolved.

14
 3.4 Dispersive power and resolving power of a grating

Figure 3.7: Diffraction spectrum of two wavelengths.

3.4.3 Resolving power of a grating

The resolving power of a diffraction grating is defined as the capacity to form
separate diffraction maxima of two wavelengths which are very close to each
other. It is given by (λ/dλ), where dλ is the difference between two wavelengths
and λ is their mean wavelength.

Figure 3.8: Diffraction of two wavelengths at a grating.

We consider a parallel beam of light consisting of two wavelengths λ and

λ + dλ as shown in Figure 3.8.
For the nth principal maximum due to λ in the direction θn (at the point P1 ),
we have

(e + d) sin θn = nλ, n = 0, ±1, ±2, ±3... (3.45)

Now, the general condition for minima is (3.32)

N (e + d) sin θ = mλ, m 6= 0, ±N, ±2N , ±3N , ..., ±nN . (3.46)

15
 3.4 Dispersive power and resolving power of a grating

Figure 3.9: Overlapping of two diffraction spectra.

The first minimum of λ adjacent to the nth principal maximum, along the
direction (θn + dθn ), will correspond to m = nN + 1, for which we get

N (e + d) sin(θn + dθn ) = (nN + 1)λ (3.47)

If the nth principal maximum of λ + dλ, occurs in the same direction (see
Figure 3.9), then

(e + d) sin(θn + dθn ) = n(λ + dλ) (3.48)

Comparing (3.47) and (3.48), we get

(nN + 1)λ = nN (λ + dλ)

=⇒ nN λ + λ = nN λ + nN dλ
λ
=⇒ = nN (3.49)
dλ

Therefore,
λ
= nN
dλ
λ N (e + d) sin θn
or, = (3.50)
dλ λ

Thus, the resolving power of a grating is directly proportional to (i) Number

of lines (slits) in the grating, and (ii) Order of the spectrum.

16
 3.5 Numerical problems - Diffraction

3.5 Numerical problems - Diffraction

Exercise 3.1
Light of wavelength 5000 Å is incident normally on a plane transmission grat-
ing. Find the difference in deviations in the first and third order spectra. No.
of lines per cm on the grating surface is 6000.

Solution:

Given, λ = 5000 Å = 5×10−5 cm.

1
= 6000/cm.
e+d
First and third order spectra, correspond to n1 = 1 and n2 = 3, respectively.

For first order spectrum,

(e + d) sin θ1 = n1 λ
n1 λ 1 × 5 × 10−5
=⇒ sin θ1 = = = 0.3
e+d 1/6000

Therefore, θ1 ' 17.460 .

For third order spectrum,

(e + d) sin θ2 = n2 λ
n2 λ 3 × 5 × 10−5
=⇒ sin θ2 = = = 0.9
e+d 1/6000

Therefore, θ2 ' 64.160 .

Therefore, difference in deviation between the two spectra = θ2 − θ1 = 46.70

Exercise 3.2
Light of wavelength 5000 Å falls normally on a plane transmission grating
having 15000 lines in 3 cm. Find the angle of diffraction for maximum intensity

17
 3.5 Numerical problems - Diffraction

in the first order.

Solution:

Given, λ = 5000 Å = 5×10−5 cm.

3 1
e+d= = cm. For first order, n = 1.
15000 5000
Therefore,

nλ 1 × 5 × 10−5
sin θ = = = 0.25
e+d 1/5000

Therefore, θ ' 14.480 .

Exercise 3.3
A single slit is illuminated by light composed of two wavelengths λ1 and λ2 . It
is observed that the first diffraction minimum for λ1 coincides with the second
diffraction minimum for λ2 . What is the relation between λ1 and λ2 ?

Solution:

For minima of single slit diffraction intensity,

e sin θ = mλ

Given, minima corresponding to m = 1 for λ1 , and m = 2 for λ2 coincide, i.e.,

occur at the same θ.
Therefore,

e sin θ = λ1 = 2λ2
=⇒ λ1 = 2λ2

Exercise 3.4
What is the minimum number of lines required in grating to just resolve the
lines of wavelength 5890 Å and 5896 Å in the second order?

Solution:

18
 3.5 Numerical problems - Diffraction

Given, λ1 = 5890 Å, λ2 = 5896 Å, n = 2.

Therefore, dλ = 6 Å and mean wavelength λ = 5893 Å.

If N is the number of lines, then
λ
nN =
dλ
λ
=⇒ N =
ndλ
×5893
= ' 491
2×6

Exercise 3.5
How many orders will be observed by a grating with 4000 lines/cm, if it is
illuminated by light of wavelength in the range 5000-7500 Å?

Solution:
For principal maxima in a grating diffraction pattern,

nλ = (e + d) sin θ

Maximum order number observable with a grating corresponds to sin θ = 1,

that is,
e+d
nmax =
λ
1
Given, (e + d) = , λ1 = 5000 Å = 5×10−5 cm, λ2 = 7500 Å = 7.5×10−5
cm. 4000
So,
e+d 1
n1max = = =5
λ1 4000 × 5 × 10−5
and
e+d 1
n2max = = = 3.3
λ2 4000 × 7.5 × 10−5
n can only take integral values, so n2max = 3.
Therefore, 3 to 5 orders will be observed for the given wavelength range.

19
 Laser
PH-1007 (Physics)
Dr. Gorky Shaw

4.1 Introduction
LASER stands for Light Amplification by Stimulated Emission of Radiation.
In the classical view, the energy of an electron orbiting an atomic nucleus
is larger for orbits further from the nucleus of an atom. However, quantum
mechanical effects force electrons to take on discrete positions in orbitals. Thus,
electrons are found in specific energy levels of an atom, and can transition
between energy levels via suitable absorption or emission of energy. An idealized
atom with two electron energy levels and one electron is shown in Figure 4.1.
These processes are more conveniently represented in energy level diagrams
as shown in Figure 4.2.

Figure 4.1: Schematic representation of the absorption, spontaneous emission

and stimulated emission processes in an idealized atom with a single
electron and two available energy levels.

4.1.1 Types of light-matter interaction

Laser action is made possible by light-matter interaction. There are three types
of possible interactions between matter (a system of atoms) and light that are
of interest:

1. Absorption is the process of moving an electron from a lower energy

level to an upper one by utilizing the energy of an incident photon (as
shown in the first panel of Figure 4.2).
Rate of absorption is proportional to the number of atoms N1 having
their electrons in the lower energy state E1 , and the number of photons
incident at a given time, n.
 4.1 Introduction

Figure 4.2: Schematic energy level diagrams showing the absorption, sponta-
neous emission and stimulated emission processes in an idealized
atom with a single electron and two available energy levels.

2. Spontaneous emission is the process in which an excited atom emits a

photon itself without any external trigger (as shown in the second panel
of Figure 4.2).
Rate of spontaneous emission is proportional to the number of atoms N2
having their electrons in the upper energy level E2 .
3. Stimulated emission is the process of emission of radiation from the
excited atom under the influence of an external trigger (such as an incident
photon, as shown in the third panel of Figure 4.2).
Rate of stimulated emission is proportional to the number of atoms N2
having their electrons in the upper energy level E2 , and the number of
photons incident at a given time, n.
At equilibrium, there is no net transition between the two energy levels,
i.e., the rate of absorption equals the rate of emission (including both
spontaneous and stimulated emission).

Some crucial differences between spontaneous emission and stimulated emis-

sion are:

1. Spontaneous emission is probabilistic in nature while stimulated emission

can be controlled.

2. In spontaneous emission, emitted light is incoherent. In stimulated emis-

sion, emitted light is coherent and monochromatic.

3. In spontaneous emission, emitted photons do not flow exactly in the same

direction as the incident photons. In stimulated emission, all emitted
photons travel in the same direction as the incident photons.

2
 4.1 Introduction

4.1.2 Some important Laser terms

1. The ground state or ground level is the lowest available energy level
in an atom. Under normal equilibrium conditions, it is the most densely
populated energy level.

2. An excited state is an upper energy level with very short lifetime. That
is, electrons reaching these energy levels from the ground state by absorp-
tion of energy (such as from incident photons) almost instantaneously
drop back down to a lower energy level with emission of energy (mostly
in the form of radiation).

3. A Metastable state is a relatively long-lived upper energy level from

where excited atoms do not return to a lower level instantaneously.
Lifetime of a highly excited state ∼ 10−8 seconds.
Lifetime of a metastable state ∼ 10−6 − 10−3 seconds.
(Lifetime of an energy level is the amount of time for which an electron
typically stays in that level before transitioning to another energy level.)

4. Population inversion is the process of achieving larger population in

a higher energy level as compared to a lower energy level. Population
inversion is used for light amplification, and is crucial for laser operation.
Population inversion cannot be achieved in a two-level system as lifetime
of the higher level is generally very short and excited atoms quickly relax
back to the ground level. An intermediate metastable level is required to
achieve population inversion.
Population inversion is schematically represented in Figure 4.3.
Let N1 , N2 , andN3 be the populations of the energy levels E1 , E2 , andE3 ,
such that E1 < E2 < E3 , as shown in Figure 4.3. In normal conditions of
thermal equilibrium, according to Boltzmann distribution law,

N1 ∝ e−E1 /kT
N2 ∝ e−E2 /kT
N3 ∝ e−E3 /kT ...etc. (4.1)

Therefore, comparing the above expressions, we get

N2
= e−[(E2 −E1 )/kT ]
N1
N3
= e−[(E3 −E2 )/kT ] ...etc. (4.2)
N2

3
 4.1 Introduction

Usually, N1  N2 > N3 (Figure 4.3(a)). Population inversion is said to

be achieved between the levels E2 and E3 when we have N3 > N2 (Figure
4.3(b)).

Figure 4.3: Schematic representation of population inversion.

5. Pumping is the method of raising the atoms from lower energy levels to
higher ones, in order to achieve population inversion.
Some pumping techniques are:
Optical pumping: Optical energy in the form of photons is used to
excite the atoms in the medium. An external light source (like xenon flash
lamp) is used to achieve population inversion the laser medium. Optical
pumping is used in solid-state lasers such as ruby lasers, discussed in
Section 4.6.
Electrical pumping: High-voltage electric discharge acts as the pump
source and direct electron excitation occurs through electric discharge.
Electrical pumping is used in gas lasers (e.g., Argon ion lasers).
Inelastic collision between atoms: In this case also, high-voltage
electric discharge acts as a pump source, but a combination of two types
of gases, with comparable excited states, is used. This method is used in
the Helium–Neon (He-Ne) laser, and described in detail in Section 4.7.
Direct conversion: In this method, electrical energy applied to semi-
conductors and LEDs is directly converted into light energy due to recom-
bination of electrons and holes. This method of pumping is used in semi-
conductor laser such as GaAs (gallium arsenide), AlGaAs (Aluminium
gallium arsenide), GaN (gallium nitride) etc.
Chemical reaction: An atom or a molecule produced to be in the
excited state at the time of production through some chemical reaction
can be used for pumping. It is used in chemical lasers such as hydrogen
fluoride (HF) laser.

4
 4.2 Three-level Laser

4.2 Three-level Laser

Figure 4.4: Energy level diagram of a three-level laser.

Consider a system with three energy levels E1 < E2 < E3 , as shown in Figure
4.4. E1 is the ground state. E3 is called the pump level. E2 is a metastable state
with lifetime much longer than that of E3 . In normal conditions, populations
in these levels are N1  N2 > N3 .
When supplied with light of energy hν = E3 − E1 , the atoms in the lower
energy state E1 are excited to the pump level E3 . This process of supplying
energy is called pumping. Out of these excited atoms, some reach the ground
state by spontaneous emission. But many of them undergo fast spontaneous
non-radiative transition to the metastable state E2 instead.
Due to the longer lifetime of the level E2 , spontaneous transitions from E2
to E1 do not occur often, leading to accumulation of atoms in the level E2 .
Eventually, we have a situation where N2 > N1 > N3 . Thus in the three-level
system, population inversion is achieved between the energy levels E1 and E2 .
A photon spontaneously emitted when an atom transitions from E2 to E1
interacts with more atoms in the level E2 and causes stimulated emission of
more photons. Photons emitted during these stimulated emission processes
again interact with atoms in the level E2 . This triggers an avalanche effect.
This way, a large number of photons are emitted and laser beam is obtained.

4.3 Four-level Laser

In this system, the energy levels are E1 < E2 < E3 < E4 as shown in Figure
4.5. Corresponding populations N1  N2 > N3 > N4 in normal thermal
equilibrium.

5
 4.3 Four-level Laser

Figure 4.5: Energy level diagram of a four-level laser.

The level E2 is above E1 by energy more than kT . Hence population of E2

is negligible compared to that of E1 at thermal equilibrium (see (4.2)), i.e., it
is almost vacant.
Similar to the three-level laser mechanism, the pump energy elevates atoms
to the short-lived pump level E4 from where atoms drop spontaneously to the
metastable level E3 . Lifetime of E3 is much longer than that of E4 and E2 .
Since the terminal level E2 is almost vacant, population inversion between
E2 and E3 is quickly established.
After completion of their lifetime in level E3 , or by an impetus of energy
hν = E3 − E2 , atoms fall to the level E2 , from where they quickly reach the
ground level E1 by radiative or non-radiative transition and are once again
available for excitation by pumping.

Advantage of a four-level system: In a a four-level system, population

inversion has to be achieved between the metastable level and a lower level
which is almost vacant to begin with. On the other hand, in a three-level
system, population inversion has to be achieved between the metastable level
and the ground level which is densely populated to begin with. Hence, in the
latter case, at least half of the population of atoms in the ground level must
be excited to higher energy levels to achieve population inversion. Thus, much
more pumping energy is required in a three-level than that in a four-level system

In both three-level and four-level systems, frequency or energy of pumping

(incident) photons must be higher than that of the emitted photons.

6
 4.4 Components of a laser system

4.4 Components of a laser system

The three principal components of a laser system (as shown in Figure 4.6) are:
1. pump source
2. laser medium
3. optical resonator.

Figure 4.6: Schematic representation of components of a laser system.

The pump source provides energy to the laser medium to attain population
inversion.
The laser medium or active medium is a medium which, when excited,
attains the state of population inversion and causes light amplification. Laser
medium can be solid, liquid, or gas.
Optical resonator: The laser medium is surrounded by two parallel mir-
rors which provide feedback of light. One mirror is perfectly reflective (high
reflector) while the other one is partially reflective (output coupler). Together
they are known as optical resonator or optical cavity or resonating cavity. The
optical coupler allows some light to leave the optical cavity to produce the laser
output beam. Photons in the laser medium bounce back and forth between the
two mirrors, causing an avalanche of stimulated emissions to achieve optical
gain.
Thus, amplified light is produced by stimulated emission and hence will travel
to large distances without spreading out in space, after having escaped through
the output coupler.

7
 4.5 Types of Laser

4.5 Types of Laser

Based on the type of laser medium used, laser systems can be classified as
follows:

1. Solid-state laser: Most prominent example is the Ruby laser. Other

host mediums include Sapphire (Al2 O3 ), Nd:YAG ((neodymium-doped
yttrium aluminium garnet), Nd:glass, Yt:glass.

2. Gas laser: He-Ne, Argon ion, CO2 , excimer, nitrogen, hydrogen.

3. Liquid laser: Dye laser.

4. Semiconductor laser: GaAs, AlGaAs, GaN.

4.6 Ruby Laser

4.6.1 Principle
Ruby laser is the first laser developed. It is a solid state three-level laser. Ruby
laser consists of a synthetic ruby crystal, Al2 O3 , doped with about 0.05% by
weight with Cr3+ ions. These ions have a set of three energy levels suitable for
the laser action.

4.6.2 Construction and working

As shown in Figure 4.7, the ruby rod is illuminated by an intense pulse of
light, which is generated by a helical Xenon discharge lamp (Xenon flash tube).
The ends of the ruby rod are highly polished and silvered to serve as laser
mirrors. The impulse of light creates the required population inversion of atoms
in the ruby rod, and due to the presence of the mirrors, the laser generation is
triggered.
When the flash tube is switched on, there is an intense burst of light lasting
for a few milliseconds. Absorption of this light energy excites many of the Cr3+
ions from the E1 energy level to the pump E3 and E4 (as shown in Figure 4.8).
The excited atoms release a part of their energy to the crystal lattice by collision
and decay to the metastable state E2 by non-radiative transitions. The level
E2 is ∼ 1.79 eV above the ground state E1 . The number of ions in the state
E2 goes on increasing with continuous pumping. Thus population inversion is
achieved between the levels E2 and E1 . During transition of an atom from E2
to E1 , a photon of wavelength 6943 Å is emitted.
The emitted photon will move back and forth between two mirrors until it
stimulates other excited ions and causes them to emit light. Thus, optical gain
is achieved.

8
 4.6 Ruby Laser

Figure 4.7: Schematic representation of the construction of Ruby laser.

The light or photons emitted due to stimulated emission will escape through
the partially reflecting mirror or output coupler to produce laser light.

Figure 4.8: Schematic representation of the energy levels and transitions in a

Ruby laser.

The ruby laser works in pulsed mode due to the high pump energy required
to operate a three-level laser. At room temperature, a ruby laser will only emit
short bursts of laser light, each laser pulse occurring after a flash of the pumping
light.

9
 4.7 He-Ne laser

It would be better to have a laser that emits light continuously. Such a laser
is called a continuous wave (CW) laser. Four-level lasers are suitable to work
in CW mode.

4.7 He-Ne laser

4.7.1 Principle
The Helium-Neon laser (also known as He-Ne laser) is the first continuous wave
(CW) laser invented. It is a type of gas laser in which a mixture of helium and
neon gas is used as a laser medium. In He-Ne lasers, electrical pumping method,
rather than optical pumping method, is used. The excitation of electrons in
the He-Ne gas active medium is achieved by passing an electric current through
the gas. The He-Ne laser operates at a wavelength of 632.8 nm.

Figure 4.9: Schematic representation of the construction of He-Ne laser.

4.7.2 Construction and working

A high voltage DC power supply, which produces energetic electrons that travel
through the gas mixture, is used as the pump source (as shown in Figure 4.9).
When the power is switched on, a high voltage of about 10 kV is applied across
the gas mixture. This power is enough to excite the electrons in the gas mixture.
The electrons produced in the process of discharge are accelerated between the
electrodes (cathode and anode) through the gas mixture. In the process of
flowing through the gas, the energetic electrons transfer some of their energy
to the helium atoms in the gas. As a result, the lower energy state electrons

10
 why He emission not possible 4.7 He-Ne laser

Figure 4.10: Schematic representation of the energy levels and transitions in a

He-Ne laser.

of the helium atoms gain enough energy and jump into the excited states or
metastable states. The metastable state electrons of the helium atoms cannot
return to ground state by spontaneous emission. However, they can return to
ground state by transferring their energy to the lower energy state electrons
of the neon atoms. The metastable states of neon have the longer lifetime.
Therefore, a large number of neon atoms will remain in the metastable states
and hence population inversion is achieved.Thus, helium atoms help neon atoms
in achieving population inversion.
Figure 4.10 shows the energy level diagram of Helium-Neon laser, with the
possible transitions. The amount of Helium in the tube is about 6 times the
amount of Neon. Thus Helium atoms have more chance to receive energy from
the accelerated electrons, and transfer into the excited energy levels E3 and E5 .
Neon atoms have two excited energy levels (E3 and E5 ) which are very close
to the excited energy levels of the Helium atom. The excited Helium atoms
transfer their excitation energy to the Neon atoms by collisions by a process
called resonance excitation.
The metastable state electrons (E3 and E5 ) of the neon atoms will sponta-
neously fall into the next lower energy states (E2 and E4 ) by releasing photons
corresponding to red light and infrared radiation. The electrons continue on to
the ground state E1 through radiative and nonradiative transitions.
The photons emitted from the neon atoms will move back and forth between
two mirrors until they stimulate other excited electrons of the neon atoms and
causes them to emit light. Thus, optical gain is achieved.
The light or photons emitted due to stimulated emission will escape through
the partially reflecting mirror or output coupler to produce laser light.
Energy from the He-Ne laser is emitted at wavelengths which correspond to

11
 4.8 Applications of Laser

the energy difference between the levels:

E5 − E4 ⇒ λ1 = 3.391 µm (near-infrared)
E5 − E2 ⇒ λ2 = 0.632 µm (visible light, red color)
E3 − E2 ⇒ λ3 = 1.152 µm (mid-infrared)

4.8 Applications of Laser

Applications of Laser include, but are not limited to:

1. Optical fibre communication

2. Cutting and welding of metals

3. Surgery

4. Optical data storage (CD, DVD etc.)

5. Laser printers, copiers etc.

6. Thermonuclear reactions

7. Holography

8. Isotope separation.

12
 Quantum Mechanics - I
PH-1007 (Physics)
Dr. Gorky Shaw

5.1 Introduction
As discussed in earlier chapters, light is a transverse, electromagnetic wave that
is visible to the human eye. Light is dualistic in nature. That is, certain optical
phenomena, such as interference, diffraction, and polarization, are explained by
the wave nature of light. Light-matter interaction in processes of emission and
absorption, such as photoelectric effect and Compton effect, are explained by
the corpuscular (particle) nature of light.
Development of quantum mechanics began in early 20th century with path-
breaking ideas of Max Planck and Niels Bohr. One of the major driving forces
towards this development was the failure of classical physics to satisfactorily
explain this wave-particle duality of light.
In particular, some of the phenomena which could not be explained by clas-
sical theories are:
1. Blackbody radiation
2. Stability of the atom
3. Photoelectric effect
4. Compton effect
5. Specific heat of solids at low temperatures

5.2 Blackbody radiation

A blackbody is an ideal body that absorbs all incident radiation (that is, there
is no reflection or transmission).
It emits (radiates) energy which is characteristic of the blackbody, and not
dependent upon the type of incident radiation. It is called blackbody ra-
diation. At uniform temperature, blackbody radiation has a characteristic
frequency (and corresponding wavelength) distribution, depending only on its
temperature. Thus, a blackbody is a hypothetical object which is a perfect
absorber and a perfect emitter of radiation over all wavelengths.
The spectral distribution of the thermal energy radiated by a blackbody
(i.e. the pattern of the intensity of radiation over a range of wavelengths or
frequencies) depends only on its temperature. The characteristics of blackbody
radiation can be described in terms of several laws, discussed below.
 5.2 Blackbody radiation
5.2.1 Stephan-Boltzmann law of blackbody radiation
Total radiant heat energy emitted from the surface of a blackbody is propor-
tional to the fourth power of its absolute temperature.

E = σT 4 (5.1)

where σ is called the Stephan-Boltzmann constant = 5.670367×10−8 Wm−2 K−4 .

5.2.2 Wien’s displacement law

The intensity of radiation emitted from a blackbody reaches its maximum at a
wavelength λm which is inversely proportional to its absolute temperature (cf.
Figure 5.1).
b
λm = (5.2)
T
where b is called the Wien’s displacement constant = 2.898 × 10−3 m · K.

Figure 5.1: Blackbody radiation as a function of wavelength for various tem-

peratures (Source: HyperPhysics).

Hence, as the temperature of a blackbody increases, its radiation peaks at

wavelengths in the visible range. Color of a blackbody at various temperatures
from 800 K to 12200 K is represented in Figure 5.2.

2
 5.2 Blackbody radiation

Figure 5.2: Color of a blackbody at various temperature from 800 K to 12200

K (Source: Wikipedia).

Stars can be approximately treated as blackbody radiators. Their visible

color depends on the temperature of the radiator.

Figure 5.3: Radiation curves for stars at different temperatures (Source: Hy-
perPhysics).

Wien’s displacement law and Stephan-Boltzmann law are in good agreement

with experimental data. However, classical laws attempting to describe the
entire spectrum of blackbody radiation (curves as shown in Figures 5.2 and
5.3) were not successful.

5.2.3 Wien’s distribution law, Rayleigh-Jeans law and the ultraviolet

catastrophe
According to Wien’s distribution law, derived from classical thermodynamic
considerations, the energy density in the wavelength range λ to (λ + dλ) at

3
 5.2 Blackbody radiation

temperature T is given by
a
Eλ dλ = dλ (5.3)
λ5 eb/λT
where a and b are constants.
According to Rayleigh-Jeans law, based on classical physical arguments and
empirical facts, the energy density in the wavelength range λ to (λ + dλ) at
temperature T is given by
8πkT
Eλ dλ = dλ (5.4)
λ4
where k is the Boltzmann constant.
As schematically represented in Figure 5.4, both of these distribution laws fail
to satisfactorily the experimental observations. Wien’s law fails to accurately
fit the data for long wavelengths. Rayleigh–Jeans law fails to describe the short
wavelength spectrum of blackbody radiation, which is also known as ultravi-
olet catastrophe as the law begins to diverge with empirical observations as
these frequencies reach the ultraviolet region of the electromagnetic spectrum.

Figure 5.4: Failure of classical theories to describe blackbody radiation.

5.2.4 Planck’s hypothesis and Planck’s law of radiation

The deficiencies were eliminated by the new radiation law proposed by Max
Planck based on his ‘quantum hypothesis’which consisted of the following main
points:

4
 5.2 Blackbody radiation

1. Source of radiations is atoms in a state of oscillation.

2. Frequency of these atomic oscillations is equal to the frequency of radia-
tion. These frequencies can have discrete values only, which are multiples
of a certain fundamental frequency (say ν)1 .
3. Oscillators can absorb or emit energy in packets of hν and switch between
energy levels with energy
nhc
En = nhν = n = 0, 1, 2, 3... (5.5)
λ
where c is the speed of light and h = 6.626 × 10−34 J·s is the Planck’s
constant.

According Planck’s radiation law, the energy density in the wavelength range
λ to (λ + dλ) at temperature T is given by
8πhc
Eλ dλ =
dλ (5.6)
λ5 ehc/λkT − 1

This law agrees very well with experimental results.

5.2.5 Classical laws from Planck’s law of radiation

For large λ, Planck’s reduces to Rayleigh-Jeans law. For large values of λ,
the term hc/λkT is small, and hence for the exponential, we can use binomial
expansion and keep the first two terms of the expansion, ignoring higher order
terms. Therefore, in this limit, (5.6) may be rewritten as
8πhc
Eλ dλ '   dλ
hc
λ5 1 + −1
λkT
8πhc
=   dλ
hc
λ5
λkT
8πkT
= dλ (5.7)
λ4
which is the Rayleigh-Jeans law (5.4).

1
In essence: every Macroscopic phenomenon which appears to be continuous is microscopically quantized.

5
 5.3 Quantization and Bohr model of the atom

For small λ, Planck’s reduces to Wien’s law. For small values of λ, ehc/λkT >> 1.
Therefore, in this limit, (5.6) may be rewritten as
8πhc
Eλ dλ ' dλ
λ5 ehc/λkT
a
= dλ (5.8)
λ5 eb/λT
which is Wien’s law (5.3), with a = 8πhc and b = hc/k.

5.3 Quantization and Bohr model of the atom

The Bohr model of the atom consists of a nucleus consisting of protons and
neutrons, with electrons orbiting around it in discrete orbits. Bohr proposed
that electrons do not radiate energy as they orbit the nucleus, but exist in
states of constant energy which he called stationary states. Atoms emit or
absorb radiation only when electrons jump between discrete energy states. The
energy change is given by

En − En−1 = hν (5.9)

which is the energy associated with a quantum of light, or photon.

Figure 5.5: Bohr model of the atom.

6
 5.4 Photoelectric effect

5.4 Photoelectric effect

5.4.1 Introduction and Classical Picture
The phenomenon of ejection (emission) of electrons from the surface of met-
als when radiation (light) of suitable frequency is incident on it is known as
Photoelectric effect.
According to the classical picture of light-matter interaction, higher frequency
of light corresponds to higher energy. Electrons would be emitted for any
frequency of light, after a delay (which is shorter for higher intensities), during
which electrons absorb sufficient energy to escape from the metal surface. This
is based on the idea that light is a wave, continuously delivering energy to the
electrons. According to this picture, at high intensities, the electrons would
absorb more energy and hence, would have greater velocity.

5.4.2 Observations
1. Photoelectric effect is an instantaneous process.
2. Photoelectric current (number of electrons ejected) is proportional to the
intensity of radiation.
3. A minimum frequency (ν0 ) of light is needed to cause photoelectric emis-
sion from a metal surface. The corresponding minimum energy required
for photoelectric emission is known as the work function (φ0 = hν0 ) of
the metal.
4. Kinetic energy of photoelectrons is proportional to the frequency of the
incident radiation.

5.4.3 Einstein’s model and the photon

Einstein explained photoelectric effect on the basis of light consisting of pho-
tons. Energy of radiation energy is made up of discrete units or localized
bundles of electromagnetic energy, called photons, which travel with velocity of
light (c) and have zero rest mass. Each photon (i.e. the quantum of energy)
hc
carries energy hν = , independent of the intensity of light. The intensity at
λ
a given frequency is determined by the number of photons incident per second.
Hence increase in intensity increases the number of photons incident per second
(but individual photon energy remains unaffected).
An electron on the metal surface completely absorbs the energy of one photon
hc
and gains the energy hν = . A part of the energy is used to free the electron
λ

7
 5.5 Particle nature of Radiation

Figure 5.6: Bohr model of the atom.

from the atom. If this quantum of energy exceeds this minimum energy needed
for the electron to escape from metal surface, i.e., the work function φ0 , then
the electron is emitted from the metal and the excess energy (hν − φ0 ) appears
1
as its kinetic energy mv 2 . Hence the equation for photoelectric emission is
2

1 2
mv = hν − φ0 (5.10)
2

where φ0 = hν0 . For ν < ν0 , there is no photoelectron emission.

5.5 Particle nature of Radiation

5.5.1 Summary of the photon picture of electromagnetic radiation
1. In interaction with matter, radiation behaves as if it is made up of parti-
cles, called photons.

2. Each photon has an energy hν and speed c (in vacuum), the speed of
light.

3. All photons of radiation of a particular frequency/wavelength have same

hc
energy hν = irrespective of the intensity of radiation. Photon energy
λ
is dependent on frequency and independent of intensity of radiation.

8
 5.5 Particle nature of Radiation

4. Photons are electrically neutral and are not deflected by electric or mag-
netic field.

5.5.2 Photon energy and momentum

Particles traveling with velocity close to the speed of light (i.e., v ∼ c) are
relativistic particles and must be treated within the domain of relativistic
mass-energy-momentum relations.
For a particle of mass m, its energy is given by

2 m0 c2
E = mc = r (5.11)
v2
1− 2
c
where m0 is the rest mass of the particle.
Since a photon moves with speed of light, we have v = c. If m0 c2 = 0, then
from (5.11), E would be infinite. But energy E of the photon is finite (= hν).
Hence, the photon must have zero rest mass (i.e., m0 = 0), and E cannot be
0
determined from this expression, which takes the indeterminate form.
0
We rather use the relativistic energy-momentum relationship

E 2 = p2 c2 + m20 c4 (5.12)

For the photon m0 = 0. Therefore, (5.12) reduces to

E = pc (5.13)

where p = E/c is the momentum of the photon.

The direction of photon momentum is in the direction propagation of radia-
tion. In a light-matter interaction (such as between photons and electrons), the
total energy and momentum are conserved. However, the number of photons
may not be conserved. Photons may be absorbed or created.

5.5.3 Momentum and dual nature of radiation

hc
Energy of a photon for radiation of frequency ν is given by E = hν = .
λ
Energy of the same photon with momentum p is given by E = pc. So, combining
these two, we have

9
 5.6 De Broglie hypothesis of matter waves

h
p=
λ
h
or, λ = (5.14)
p

5.6 De Broglie hypothesis of matter waves

Louis de Broglie argued that nature is symmetric, and what is true for photons
(radiation) must be true for material particles. He suggested that wave parti-
cle duality is a universal principle and is not restricted only to electromagnetic
radiations. Hence, material particles must exhibit wave-like properties. De
Broglie proposed that all material particles can manifest themselves as waves,
called matter waves (also called de Broglie waves).

The wavelength of matter waves (the de Broglie wavelength) associated with

a moving particle of linear momentum p = mv is given by
h h
λ= = (5.15)
p mv

De Broglie proposed that the role of matter waves in guiding the motion of
material particles is analogous to that of electromagnetic waves in guiding the
motion of photons.

5.6.1 De Broglie wavelength of non-relativistic particles

For non-relativistic particles v << c and m ≈ m0 . In this case
h h h
λ= = ≈ (5.16)
p mv m0 v
Further, for a non-relativistic free particle (i.e., with not subject to a potential),
p2
the total energy E is equal to its kinetic energy K = . That is,
2m
√ √
p = 2mK = 2mE (5.17)

10
 5.6 De Broglie hypothesis of matter waves

Hence,
h h h
λ= =√ =√ (5.18)
p 2mK 2mE

5.6.2 De Broglie wavelength associated with an electron

If a particle with charge q is accelerated through a potential difference of V
volt, then kinetic Energy gained

p2
K= = qV (5.19)
2m
Hence, in this case
h h h
λ= =√ =√ (5.20)
p 2mK 2mqV

For an electron, m = me = 9.1 × 10−31 kg, q = e = 1.6 × 10−19 C. Using these

values in (5.20), the de Broglie wavelength associated with an electron, λe is
given by
12.27
λe = √ Å (5.21)
V

For example, for V = 100 V, λe = 1.227 Å.

5.6.3 De Broglie wavelength of relativistic particles

From (5.11), mass of relativistic particles (v ∼ c) varies with velocity v as
m0
m= r (5.22)
2
v
1−
c2
Hence
m0 v
p = mv = r (5.23)
v2
1− 2
c

11
 5.6 De Broglie hypothesis of matter waves

and
r
v2
h h 1 − c2
λ= = (5.24)
p m0 v

Note: If the term pc in (5.12) is negligible compared to m0 c2 , then the

motion is non-relativistic. If the two terms are comparable, then the motion is
relativistic.

5.6.4 Properties of matter waves

1. Lighter particles have longer wavelengths.

2. Higher velocity corresponds to shorter wavelength.

3. For v = 0, λ = ∞ (i.e., there is no wave). Hence, matter waves are

generated only when a particle is in motion.

4. Velocity of matter waves depends on the velocity of the particle.

5. Velocity of matter waves may be greater than the velocity of light!

(The last two points are discussed in the following sections.)

5.6.5 Matter waves and structure of the atom

For an electron to exist in a stationary state around a nucleus, the associated
electron wave must form a standing wave, i.e., the circumference must be mul-
tiples of the electron wavelength as shown in Figure 5.7.

2πrn = nλe (5.25)

Figure 5.7: Electron orbits as standing waves around the nucleus (Source:CK-12
Foundation).

12
 5.6 De Broglie hypothesis of matter waves

A related result is that the angular momentum of the electron in an orbit of

radius rn ,
hrn hrn nh
L = mvrn = = = (5.26)
λe 2πrn 2π
n

is quantized.

5.6.6 Phase velocity of de Broglie matter waves

Phase velocity (vp ) of a wave is the rate which phase of the wave propagates in
a medium.

vp = νλ (5.27)

For a simple harmonic wave propagating along the +ve x-direction, given by

y = a sin(ωt − kx) (5.28)

where a is the amplitude and (ωt − kx) is the phase of the wave, the phase
velocity
dx
vp = (5.29)
dt
is given by
2πν ω
vp = νλ = = (5.30)
2π/λ k

(5.30) is obtained as follows:

For a plane wave of constant phase,

(ωt − kx) = constant

d
=⇒ (ωt − kx) = 0
dt
dx
=⇒ ω − k =0
dt
dx ω
=⇒ vp = = (5.31)
dt k

13
 5.6 De Broglie hypothesis of matter waves

At this point, we introduce the reduced Planck’s constant

h
~= = 1.055 × 10−34 J · s (5.32)
2π

In terms of the reduced Planck’s constant, the energy and momentum of the
particle associated with the wave are given by
h
E = hν = × 2πν = ~ω
2π   
h h 2π
and, p = = · = ~k (5.33)
λ 2π λ

Accordingly, the phase velocity can also be expressed as the ratio of the energy
and momentum of the associated particle:

ω ~ω E
vp = = = (5.34)
k ~k p

5.6.7 Relationship between phase velocity and particle velocity

1. Non-relativistic particle with velocity v:

1 p2
E = mv 2 = (p = mv) (5.35)
2 2m
Therefore,

E p2 p mv v
vp = = = = = (5.36)
p 2mp 2m 2v 2
which indicates the wave travels at half the speed of the associated par-
ticle.

2. Relativistic particle with velocity v:

E = mc2 and p = mv (5.37)

14
 5.6 De Broglie hypothesis of matter waves

Therefore,

E mc2 c2
vp = = = (5.38)
p mv v
Now, for material particles, v < c always. Hence, we have vp > c.

which indicates the wave travels faster than light!

Hence, from (5.37) and (5.38), we note that for both non-relativistic and rela-
tivistic particles, the phase velocity of matter wave is different from the veloc-
ity of the associated particles. This apparent anomaly was resolved by Erwin
Schrödinger by introducing the wave packet representation of matter waves,
i.e., considering that a packet (or group) of waves, rather than a single wave, is
associated with a particle in motion.

5.6.8 Wave packet and group velocity

A wave packet (or wave group) is a resultant of a group of waves with slightly
differing velocities and wavelengths interfering constructively over a small region
of space where the particle can be located (∆x in Figure 5.8). Outside this
space, they interfere destructively so that the resultant amplitude is zero.

Figure 5.8: A wave packet.

The velocity of a wave packet associated with a moving particle is known as

the group velocity. To obtain an expression for group velocity, we consider
the simplest scenario of a wave packet arising from the combination of two

15
 5.6 De Broglie hypothesis of matter waves

Figure 5.9: Phase velocity and group velocity of a wave packet.

simple harmonic waves with the same amplitude but with a difference of dω
and dk in angular frequency and wave number, respectively, given by2

y1 = a cos(ωt − kx)
and y2 = a cos[(ω + dω)t − (k + dk)x] (5.39)

The resultant displacement

y1 = y1 + y2 = a cos(ωt − kx) + a cos[(ω + dω)t − (k + dk)x]

= a{cos(ωt − kx) + cos[(ω + dω)t − (k + dk)x]}
   
(2ω + dω)t − (2k + dk)x dωt − dkx
= 2a cos cos (5.40)
2 2

Assuming dω << ω and dk << k, we have

 
dω dk
y = 2a cos(ωt − kx) cos t− x (5.41)
2 2

2
we could as well use sin instead of cos, the end result would be the same.

16
 5.6 De Broglie hypothesis of matter waves

Reordering,
 
dω dk
y = 2a cos t − x cos(ωt − kx) (5.42)
2 2

(5.42) represents a wave of higher angular frequency ω and wave number k

dω
superimposed on a slower modulation envelop with angular frequency and
2
dk
wave number (as represented in Figure 5.9). That is, it can be rewritten as
2
y = A cos(ωt − kx) (5.43)

where
 
dω dk
A = 2a cos t− x (5.44)
2 2

is the modulated amplitude of the resultant wave.

It represents the envelope of the wave group. Group velocity, the velocity of
propagation of the envelop, and hence of the modulated wave group, is given
by
dω
vg = (5.45)
dk

5.6.9 Relationship between group velocity and particle velocity

From (5.33) and (5.45),

dω d(~ω) dE
vg = = = (5.46)
dk d(~k) dp
1. Non-relativistic particle with velocity v:
 2
d p p
vg = = =v (5.47)
dp 2m m

17
 5.6 De Broglie hypothesis of matter waves

2. Relativistic particle with velocity v:

2pc2 pc2 mvc2

q 
d 2 2 2 4
1
vg = p c + m0 c = · p = = =v
dp 2 p2 c2 + m20 c4 E mc2
(5.48)

From (5.47) and (5.48), for both non-relativistic and relativistic particles, we
have

vg = v (5.49)

That is, group velocity is always equal to particle velocity. Which mean the de
Broglie wave group always travels along with the particles, hence resolving the
anomaly associated withe the phase velocity and particle velocity.

5.6.10 Relationship between group velocity and phase velocity

From (5.34), phase velocity
ω
vp =
k
=⇒ ω = vp k (5.50)

Therefore group velocity

dω d
vg = = (vp k)
dk dk
dvp
= vp + k
dk
dvp dλ
= vp + k · (5.51)
dλ dk
2π
Now, λ = . Therefore,
k
 
dλ d 2π 2π
= =− (5.52)
dk dk k k2

Using (5.52) in (5.51), we get

dvp 2π
vg = vp − k · (5.53)
dλ k 2

18
 5.7 Heisenberg’s uncertainty principle

Or,

dvp
vg = vp − λ (5.54)
dλ

dvp
In a normal dispersive medium3 , > 0. Therefore, vg < vp .
dλ
If vp is independent of λ (i.e., in a non-dispersive medium, such as vacuum),
then vg = vp .

5.7 Heisenberg’s uncertainty principle

For a narrow wave packet, precise determination of the particle position is
possible (small ∆x), but estimation of the corresponding de Broglie wavelength
is difficult (large ∆λ). For a wide wave packet, wavelength can be estimated
more accurately, but not the particle position (as illustrated in Figure 5.10).
From (5.15), uncertainty in wavelength (∆λ) corresponds to uncertainty in
momentum (∆p).

Heisenberg’s uncertainty principle: It is impossible to simultaneously

know both the exact position and the exact momentum of an object. Mathe-
matically,

∆p · ∆x = h (5.55)

The modified uncertainty principle, based on experimental observations, is

h ~
∆p · ∆x ≥ = (5.56)
4π 2

5.7.1 Justification of Heisenberg’s uncertainty principle

As in Section 5.6.8, the resultant of a wave group arising from the combination
of two waves represented by

y1 = a cos(ωt − kx)
and y2 = a cos[(ω + ∆ω)t − (k + ∆k)x] (5.57)
3
A dispersive medium is a medium in which waves of different wavelengths travel at different velocities.

19
 5.7 Heisenberg’s uncertainty principle

Figure 5.10: Heisenberg’s uncertainty principle (Source: reddit r/askscience).

is given by
 
∆ω ∆k
y = 2a cos t− x cos(ωt − kx) (5.58)
2 2

(assuming ∆ω << ω and ∆k << k).4 Wavelength of modulation of the wave

group is given by
2π 4π
λm = = (5.59)
∆k/2 ∆k

The particle could be found in any of the regions of width λm /2, as shown in
Figure 5.11. Hence, the uncertainty in particle position is
λm 2π
∆x = = (5.60)
2 ∆k
Now since de Broglie wavelength associated with the particle with momentum

4
Note the change in notation.

20
 5.7 Heisenberg’s uncertainty principle

Figure 5.11: Uncertainty in particle position.

p is
h
λ= (5.61)
p
we have
2π 2πp
k= = (5.62)
λ h
and
2π
∆k = ∆p (5.63)
h
Therefore, from (5.60) and (5.63),
2πh h
∆x = =
2π∆p ∆p
=⇒ ∆x · ∆p = h (5.64)

In general, wave packets can have different shapes, and we have

∆x · ∆p ≥ h (5.65)

As mentioned above ((5.56)), more rigorous calculations show the formulation

of position- momentum uncertainty principle is given by
~
∆x · ∆p ≥ (5.66)
2

21
 5.7 Heisenberg’s uncertainty principle

Other form of the uncertainty principle are:

Energy - time uncertainty principle:

~
∆E · ∆t ≥ (5.67)
2
where, ∆E and ∆t are the uncertainty in the energy and time, respectively.

Angular position - angular momentum uncertainty principle:

~
∆θ · ∆L ≥ (5.68)
2
where ∆θ and ∆L are the uncertainty in the angular position and angular
momentum, respectively.

5.7.2 Application: non-existence of electrons inside the nucleus

First let us assume an electron exists inside the nucleus.
The diameter of a typical nucleus is ∼ 10−14 m. This is the uncertainty ∆x
in the position of a particle existing inside the nucleus.
Therefore, corresponding uncertainty in its momentum,
~
∆p =
2∆x
1.055 × 10−34
=
2 × 10−14
≈ 0.527 × 10−20 kg m s−1 (5.69)

Now, the momentum p must at least be of the same order as ∆p and corre-
sponding velocity should at least be of the order of the uncertainty in velocity

∆p 0.527 × 10−20
∆v = = −31
= 0.0579 × 1011 ≈ 5.8 × 109 > c (5.70)
m 9.1 × 10
which is not possible.
Further, energy is required in order to confine a particle in a given volume.
Energy of an electron inside the nucleus would be at least
q
E = p2 c2 + m20 c4 (5.71)

22
 5.7 Heisenberg’s uncertainty principle

Now,

m0 c2 = 9.1 × 10−31 × (3 × 108 )2 = 81.9 × 10−15 J ' 0.511 MeV (5.72)

and

pc = 0.527 × 10−20 × 3 × 108 = 1.581 × 10−12 J ' 10 MeV (5.73)

Therefore,

E ≥ 10 MeV (5.74)

However, experiments show that electrons emitted by certain unstable nuclei

(by processes such as beta-decay) never have more than a small fraction of this
energy. Hence electrons cannot reside inside the nucleus.

23
 Quantum Mechanics - II
PH-1007 (Physics)
Dr. Gorky Shaw

6.1 The wave function

1. The quantity whose variations make up the matter waves is known as the
wave function. It is a purely mathematical quantity.

2. The wave function associated with a particle describes its position and
state. If we know the wave function ψ(x, y, z, t) at all points in space and
all instants of time, then the particle behavior can be completely specified.

3. ψ(x, y, z) represents the stationary state of the particle, independent of

time.

4. ψ is usually a complex quantity.

5. For a complex function ψ(x, y, z, t), its complex conjugate is written as

ψ ∗ (x, y, z, t), and ψ ∗ ψ = |ψ|21 .

6. ψ 2 or |ψ|2 is a real quantity and has a physical interpretation, but not

ψ. Value of this quantity at any point (x, y, z) at an instant of time t
gives the probability density (probability per unit volume) of finding the
particle at that point.

7. Probability density ρ(x, y, z, t) = ψ ∗ ψ = |ψ|2 . This probability density so

defined is positive.

8. The wave function ψ(x, y, z, t) must be finite, single-valued, and continu-

ous over the complete range of variables.

9. Probability of finding a particle in an elemental volume dV = dxdydz is

given by
|ψ|2 dxdydz
In one-dimension, probability of finding a particle within an elemental
length dx is given by
|ψ(x, t)|2 dx

1
Given a complex number z = a + ib, a, b ∈ R, it can be shown that its complex
√ conjugate is z ∗ = a − ib.
The normal length of z is a positive (and thus real) number number, |z| = a + b2 . It can be verified that
2
2
|z| = zz ∗ .
 6.1 The wave function

Note: From this point onward, our discussion will be strictly limited to
one-dimension, unless otherwise specified.

10. Integral version: The probability Pab (t) of finding a particle in the interval
x ∈ [a, b] is,
Z b Z b
2
Pab (t) = |ψ(x, t)| dx = ρ(x, t)dx (6.1)
a a

11. Total probability of finding the particle somewhere in space, if it exits, is

unity. That is, Z ∞
|ψ(x, t)|2 dx = 1 (6.2)
−∞
A wave function satisfying (6.2) is known as a normalized wave function.
Every wave function representing matter waves must be normalized.
12. In general, a wave function may not be a normalized wave function. To
normalize a wave function is to multiply by a constant factor so that the
condition (6.2) is satisfied.
Let ψ be a wave function which is not normalized. If ψ is a solution of a
wave equation, then N ψ, where N is a complex number, is also a solution.
For N ψ to be a normalized function, we must have
Z ∞
(N ψ)∗ (N ψ)dx = 1
−∞
Z ∞
=⇒ |N | 2
ψ ∗ ψdx = 1
−∞
1
=⇒ |N |2 = R ∞ ∗
(6.3)
−∞ ψ ψdx

N is called the normalization constant and N ψ the normalized wave func-

tion.
13. The wave functions of matter waves obey the principle of superposition.
That is, if ψ1 and ψ2 are solutions of a wave equation, then
ψ = a1 ψ1 + a2 ψ2
where a1 and a2 are complex constants, is also a solution.
14. The wave function associated with a free particle represents a sinusoidal
wave, implying a precisely determined momentum and completely uncer-
tain position.

2
 6.2 Schrödinger equation

6.2 Schrödinger equation

Immediately after the publication of de Broglie hypothesis, Erwin Schrödinger
proposed that the behavior of matter waves representing a material particle
is governed by a certain differential equation for the wave function ψ. The
Schrödinger equation plays the same role in wave mechanics as Newton’s second
law or Newtonian equation in classical mechanics. We cannot rigorously derive
the Schrödinger equation from existing physical principles, rather it can be
formulated as discussed below.

6.2.1 Time-dependent Schrödinger equation (TDSE)

Let the wave associated with a (non-relativistic) particle be represented by the
wave function

ψ(x, t) = Aei(kx−ωt) (6.4)

Differentiating (6.4) w.r.t. t,

∂ψ
= −iωAei(kx−ωt) = −iωψ (6.5)
∂t
Differentiating (6.4) w.r.t. x, twice,
∂ψ
= ikAei(kx−ωt)
∂x
∂ 2ψ
=⇒ 2
= i2 k 2 Aei(kx−ωt) = −k 2 ψ (6.6)
∂x
Now, the energy E and momentum p of the particle are given by

E = hν = ~ω
h
and, p = = ~k (6.7)
λ
The energy and momentum are related as

p2
E= + V (x, t) (6.8)
2m
where V (x, t) is the potential energy of the particle.

3
 6.2 Schrödinger equation

Using (6.7) in (6.8), we get

~2 k 2
~ω = +V (6.9)
2m
Multiplying (6.9) by ψ, we get

~2 k 2
~ωψ = ψ+Vψ (6.10)
2m
Now, from (6.5),
~ ∂ψ ∂ψ
~ωψ = − = i~ (6.11)
i ∂t ∂t
From (6.6),

~2 k 2 ~2 ∂ 2 ψ
ψ=− (6.12)
2m 2m ∂x2
Therefore, using (6.11) and (6.12) in (6.10), we get

∂ψ ~2 ∂ 2 ψ
i~ =− +Vψ (6.13)
∂t 2m ∂x2
This is the time-dependent Schrödinger equation (TDSE).

6.2.2 Time-independent Schrödinger equation (TISE)

In case the energy of the system does not change with time, i.e., then E = ~ω,
from (6.11) we have
∂ψ
i~ = ~ωψ = Eψ (6.14)
∂t
Using (6.14) in (6.13), we get

~2 d2 ψ
Eψ = − +Vψ (6.15)
2m dx2
Rearranging (6.15), we get

4
 6.2 Schrödinger equation

d2 ψ 2m
+ 2 (E − V )ψ = 0 (6.16)
dx2 ~
This is the time-independent Schrödinger equation (TISE).

Note that in (6.16), the partial derivative is replaced by total derivative as

there is no time dependence.

6.2.3 Time-independent Schrödinger equation (TISE): rigorous approach

In case the potential V has no time dependence (i.e., in this case, V = V (x)
only), a stationary state of constant energy E can be obtained by separating
the position and time components of ψ(x, t).
We rewrite the wave function given in (6.4) as

ψ(x, t) = Aei(kx−ωt)
= Aeipx/~−Et/~
= Aeipx/~ e−iEt/~
= φ(x)e−iEt/~ (6.17)

where

φ(x) = Aeipx/~ (6.18)

is a function of x only.
Using (6.18) in (6.13), we get

∂ h i  ~2 ∂ 2 h i
−iEt/~ −iEt/~
i~ φ(x)e = − + V (x) φ(x)e
∂t 2m ∂x2
2 2
   h
iE −iEt/~ ~ ∂ −iEt/~
i
=⇒ i~ − φ(x)e = − + V (x) φ(x)e
~ 2m ∂x2
(6.19)

Canceling out the exponentials and simplifying (6.19), we get

~2 ∂ 2
 
Eφ(x) = − + V (x) φ(x) (6.20)
2m ∂x2

Since there is no longer any time dependence in (6.20), we replace the partial
derivative with total derivative. Hence, we get

5
 6.3 Observables, operators, eigenfunctions, eigenvalues, expectation values

~2 d2
 
Eφ(x) = − + V (x) φ(x)
2m dx2
~2 d2
 
or simply, Eφ = − +V φ (6.21)
2m dx2

Rearranging (6.21), we get

d2 φ 2m
+ 2 (E − V )φ = 0 (6.22)
dx2 ~
This is the time-independent Schrödinger equation (TISE).

In subsequent Sections, we discuss applications of the Schrödinger equation.

6.3 Observables, operators, eigenfunctions, eigenvalues,

expectation values
An observable is a quantity obtained by the process of observation or mea-
surement on a physical system.
It is always a real number as it is the result of a measurement.

6.3.1 Assumptions in quantum mechanics

1. An observable depends on the measurement even though it is a charac-
teristic of the system.

2. Some observables have some interpretation as in classical mechanics (such

as momentum, total energy, position etc.) while others don’t (such as
spin).

3. Simultaneous measurement of conjugate variables (position-momentum,

time-energy, etc.) is governed by the uncertainty principle.

Associated with each observable in a physical system is a quantum mechanical

operator. Z
d
In general, an operator O is a mathematical operation (such as +, −, ×, ÷, , dx
dx
etc.) performed on a function f (x) that outputs another function g(x).

6
 6.3 Observables, operators, eigenfunctions, eigenvalues, expectation values

Mathematically,

Of (x) = g(x) (6.23)

For example,
d 2
(x ) = 2x (6.24)
dx
Order of operation is important. For example,
d d d
[xf (x)] = f (x) + x f (x) 6= x f (x) (6.25)
dx dx dx
That is, operators are not necessarily commutative.
In quantum mechanics, an operator applied to a wave function gives the
corresponding observable quantity, multiplied by the wave function.
For example, the Hamiltonian operator Ĥ

Ĥψ = Eψ (6.26)

gives the total energy E of the system.

6.3.2 Examples of operators and corresponding observables

1. Cartesian coordinate (or position) operator. Represented by x̂, ŷ, ẑ. Cor-
responding observables are x, y, z.

x̂ψ = xψ
ŷψ = yψ
ẑψ = zψ (6.27)

2. Any other function of position, e.g., potential energy V̂ (x).

V̂ (x)ψ = V (x)ψ (6.28)

3. Linear momentum operator

∂
pˆx = −i~ (x-component)
∂x
p̂ = −i~∇ (3-D) (6.29)

4. Hamiltonian or total energy operator

7
 6.3 Observables, operators, eigenfunctions, eigenvalues, expectation values

(a) Time-dependent:
∂
Ê(x) = i~ (6.30)
∂t

(a) Time-independent (Hamiltonian):

~2 ∂ 2
 
Ĥ(x) = − + V (x) (6.31)
2m ∂x2

Using this in (6.21), we get

Eφ = Ĥφ
or, Ĥψ = Eψ (6.32)

which is a compact (and the most famous) form of the time-independent

Schrödinger equation2 .

5. Kinetic energy operator

~2 ∂ 2
K̂(x) = − (x-component)
2m ∂x2
~2 2
K̂ = − ∇ (3-D) (6.33)
2m

6.3.3 Eigenvalues and eigenfunctions

A function ψ(x) is an eigenfunction of an operator  is application of  on
ψ(x) gives ψ(x) again, times a constant.

Âψ = aψ (6.34)

The possible values of the constant a are called eigenvalues of the observable
associated with the operator Â. The corresponding state of the system is known
as the eigenstate of the associated observable.
d2
For example, Let ψ(x) = sin 2x and  = 2 .
dx
d2 d
Âψ = (sin 2x) = (2 cos 2x) = −4 sin 2x = −4ψ (6.35)
dx2 dx

2
Note that the wave function symbol is changed from φ to ψ in (6.32). This of course does not change the
result.

8
 6.3 Observables, operators, eigenfunctions, eigenvalues, expectation values

d2
Hence, sin 2x is an eigenfunction and −4 an eigenvalue of the operator .
dx2
d
On the other hand, Let ψ(x) = sin 2x and  = .
dx
d
Âψ = (sin 2x) = 2 cos 2x 6= aψ (6.36)
dx
d
Hence sin 2x is not an eigenfunction of the operator .
dx

6.3.4 Expectation value

Expectation value hAi of an observable A is given by
R ∗
ψ Âψdx
hAi = R ∗ (6.37)
ψ ψdx

For a normalized wave function, ψ ∗ ψdx = 1. In this case,

R

Z
hAi = ψ ∗ Âψdx (6.38)

where the integral is over all available space and hAi may vary in time.

For example, for position x (in 1-D), the expectation value

Z Z
hxi = ψ ∗ x̂ψdx = ψ ∗ xψdx (6.39)

When a system is in an eigenstate of an observable A (i.e., when the wave

function is an eigenfunction of Â), the expectation value of A is the eigenvalue
of the wave function.
That is, if Âψ = aψ, then
Z
hAi = ψ ∗ Âψdx
Z
= ψ ∗ aψdx
Z
= a ψ ∗ ψdx
=a (6.40)

9
 6.4 Particle in an one-dimensional potential well of infinite height

That is,

hAi = a (6.41)

6.4 Particle in an one-dimensional potential well of infinite

height
As shown in Figure 6.1, we consider a particle of mass m moving along the
x−axis in a one-dimensional box of infinite height, with potential energy of the
box defined as

V (x) = ∞ for x < 0 and x > L

V (x) = 0 for 0 ≤ x ≤ L (6.42)

The walls at x = 0 and x = L are considered rigid, i.e., there is no loss of

Figure 6.1: Particle in an one-dimensional potential well of infinite height

(Source: Chemistry LibreTexts).

particle energy on bouncing back from the walls.

We consider the Schrödinger equation (6.22)

d2 ψ 2m
+ 2 (E − V )ψ = 0 (6.43)
dx2 ~
Now inside the box, V = 0. Therefore, Schrödinger equation of the particle

10
 6.4 Particle in an one-dimensional potential well of infinite height

inside the box is given by

d2 ψ 2mE
+ 2 ψ=0 (6.44)
dx2 ~
Now,

p2 ~2 k 2 ~2 k 2
E= +V = +0=
2m 2m 2m
2mE
=⇒ 2
= k2 (6.45)
~
Using (6.45) in (6.44), we get

d2 ψ
2
+ k2ψ = 0 (6.46)
dx
General solution of (6.46) can be written as

ψ(x) = A1 eikx + A2 e−ikx

or, ψ(x) = A sin(kx) + B cos(kx) (6.47)

where A1 , A2 , A and B are constants that can be determined from appropriate

boundary conditions.
Now, the particle cannot penetrate the rigid walls, or exist outside the box.
Hence we have the boundary condition

ψ = 0 at x = 0 (6.48)

and

ψ = 0 at x = L (6.49)

Using (6.48) in (6.47), we get

0 = A sin 0 + B cos 0
=⇒ 0 = 0 + B
=⇒ B = 0 (6.50)

Therefore,

ψ(x) = A sin(kx) (6.51)

11
 6.4 Particle in an one-dimensional potential well of infinite height

Using (6.49) in (6.51), we get

0 = A sin(kL)
=⇒ kL = nπ , n = 0, 1, 2, ... (6.52)

(the other possibility from (6.52), i.e., A = 0 is discarded because that would
imply ψ = 0 everywhere, which means the particle would not exist at all).
Further, n = 0 =⇒ k = 0, which once again means ψ = 0 for all x, and the
particle would not exist. Hence n = 0 is also discarded. Therefore, we have

kL = nπ , n = 1, 2, ...
nπ
=⇒ k = (6.53)
L
Therefore, we have
 nπx 
ψn (x) = A sin (6.54)
L
where we have added the suffix n to ψ to indicate dependence of ψn on n. Now,
the wave function must be normalized. That is,
Z L
|ψ(x, t)|2 dx = 1
0
Z L
2 nπx
 
2
=⇒ A sin dx = 1
L
2 Z L
0  
A 2nπx
=⇒ 1 − cos dx = 1
2 0 L
A2
  
L 2nπx
=⇒ x− sin =1
2 2nπ L
A2
  
L 2nπL
=⇒ L− sin =1
2 2nπ L
A2 L
 
1
=⇒ 1− sin (2nπ) = 1
2 2nπ
A2 L
=⇒ =1
2 r
2
=⇒ A = (6.55)
L
Therefore, we have

12
 6.4 Particle in an one-dimensional potential well of infinite height

r
2  nπx 
ψn (x) = sin (6.56)
L L
which is the wave function associated with a particle enclosed in an infinitely
deep potential well.

The wave function given by (6.56) can be written in terms of the wavelength
2L
λn = :
n
r  
2 2πx
ψn (x) = sin (6.57)
L λn

The corresponding energies are given by

~2 k 2
En =
2m
~2 n2 π 2 n2 π 2 ~2
= · 2 = , n = 1, 2, 3... (6.58)
2m L 2mL2

Therefore, energy levels of a particle inside an infinitely deep potential well

are quantized and given by

n2 π 2 ~2 n2 h2
En = = , n = 1, 2, 3... (6.59)
2mL2 8mL2

Hence, we have

En ∝ n2
and, En = n2 E1 (6.60)

Note that the lowest energy level is not E = 0. Which means a particle in a
box has a zero-point energy 3 .
The variation of wave function (ψn ), probability density (|ψn |2 ), and energy
levels (En ) of the particle are shown in Figure 6.2.
A significant feature of the states of the particle in a box is the occurrence of
nodes. These are points, other than the two end-points (which are fixed by the

3
The zero-point energy is fundamentally related to the Heisenberg uncertainty principle.

13
 6.5 Quantum tunneling of particles through potential barriers

boundary conditions), at which the wave function vanishes. At a node there is

exactly zero probability of finding the particle. The nth energy eigenstate has,
in fact, (n − 1) nodes4 .

Figure 6.2: (a) Wave function, (b) Probability density, and (c) Energy levels of
a particle in an one-dimensional potential well of infinite height.

6.5 Quantum tunneling of particles through potential

barriers
Quantum tunneling is a phenomenon in which particles penetrate a potential
energy barrier with height greater than the total energy of the particles.
A classical analogy: Consider a ball with energy E encounters a hill while
rolling. If the height of the hill corresponds to a potential energy U0 < E, then
the ball will roll over. But if the U0 > E, the ball would reach some height,
then return in the direction it came from. It would never make it to the other
side in this case.
However, quantum mechanically, the ball has a wave function defined over all
space. There is always a probability that as the ball encounters the hill, it will
“tunnel”through the hill and suddenly be found on the other side. The wave
function may be highly localized, but this probability is appreciable if the wave
packet is wider than the barrier.
A shown in Figure 6.3(a), we define a finite height rectangular potential

4
Also, easy to remember: n gives the number of crests that the wave function has inside the box.

14
 6.5 Quantum tunneling of particles through potential barriers

Figure 6.3: Quantum tunneling across a barrier of finite height.

barrier by the potential energy function

U (x) = 0 for x < 0 and x > L

U (x) = U0 for 0 ≤ x ≤ L (6.61)

If U0 → ∞, then no tunneling is possible, just as in the classical case.

We denote the wave functions in regions I, II and III as ψI , ψII , ψIII , respec-
tively, as shown Figure 6.3(b). Time-independent Schrödinger equation in each
region is:

d2 ψI 2mE
Region I: + 2 ψI = 0 (6.62)
dx2 ~

d2 ψII 2m
Region II: + 2 (E − U0 )ψII = 0 (6.63)
dx2 ~

d2 ψIII 2mE
Region III: + 2 ψIII = 0 (6.64)
dx2 ~
Evidently, ψI and ψIII will be sinusoidal functions and ψII will be an exponential
function 5 , as qualitatively represented in Figure 6.3(b).
5
This should be obvious from Section 6.4 and earlier discussion on damped oscillations (solution for overdamped
oscillations. Try to figure it out!)

15
 6.5 Quantum tunneling of particles through potential barriers

6.5.1 Footnote: Quantum mechanical tunneling: myth or reality

Quantum mechanical tunneling sounds rather exotic, but is indeed very real,
and is the underlying principle which makes a lot of things work, including, but
not limited to:

1. Josephson junctions - Electron pairs tunnel through an insulating layer

separating an two superconducting layers.

2. Radioactive decay - α-particles having kinetic energy of a few MeV (4-9)

are able to escape from the nucleus whose potential wall us typically ∼ 25
MeV high.

3. Tunnel diodes - Electrons tunnel through a potential barrier (at the con-
tact point of two different semiconductors) much larger than their kinetic
energies.

Figure 6.4: Schematic diagram of a scanning tunneling microscope (Source:

Chemistry LibreTexts).

4. Scanning tunneling microscopy (STM) - This technique allows imaging of

individual atoms on the surface of a metal. When a STM tip is brought

16
 6.5 Quantum tunneling of particles through potential barriers

close to conducting surface with a voltage bias, electrons are able to tunnel
between the tip and the sample, generating an electric current.
By using piezoelectric regulators, the tip-sample height can be adjusted
to keep the tunneling current constant. Thus, path of the tip along the
sample surface can be used to obtain a gray-scale image representing the
topography of the sample surface.

********************* The End (of this Unit) *********************

17
 Vector Algebra
PH-1007 (Physics)
Dr. Gorky Shaw

6.1 Some definitions

1. Scalars or scalar functions are quantities fully described by magnitude
only.

2. Vectors or vector functions are quantities fully described by both mag-

nitude and direction.

3. Point function is a continuous function of position of a point in a space.

4. Field is a region of space where a point function specifies a physical

quantity.

5. Scalar field is a region where, at each point, a scalar function has a

definite value.

6. Vector field is a region where, at each point, a vector function has a

definite value.

6.2 Vector products

6.2.1 Vector scalar product or dot product

Figure 6.1: Dot product of two vectors.

~ and B
The scalar product (or dot product) of two vectors A ~ is defined by

~·B
A ~ = AB cos ϕ (6.1)
 6.2 Vector products
where A, B are magnitudes (lengths) of the vectors A, ~ B, ~ respectively, and ϕ
is the angle between the vectors, as shown in Figure 6.1(a). The dot product
between of vectors is based on the projection of one vector in the direction of
the other, as illustrated in Figures 6.1(b) and 6.1(c). The result of the dot
product itself is a scalar, hence the name scalar product.
The scalar product of two orthogonal vectors (ϕ = 900 ) vanishes. In this case

~·B
A ~ = AB cos 900 = 0 (6.2)

The scalar product of a vector with itself gives the square of its magnitude:
~·A
A ~ = A2 cos 00 = A2 (6.3)

6.2.2 Vector product or vector cross product

Figure 6.2: Cross product of two vectors.

~ and B
The cross product of two vectors A ~ is defined by

~ =A
C ~×B
~ = AB sin ϕ n̂ (6.4)

The result of the cross product is a vector along a direction perpendicular to

both A~ and B,~ as indicated in Figure 6.2(a). n̂ is a unit vector (defined in
Section 6.3).
The direction perpendicular to the plane that contains the vectors A ~ and B
~
could be either ‘up’ or ‘down’, and the direction of the vector product may be

2
 6.3 Vector algebra - component form

either one of them. The direction of the vector product is determined by the
direction of the angle ϕ, as indicated in Figures 6.2(a) and 6.2(b).
As evident from Figure 6.2,
~ × B)
(A ~ = −(B
~ × A)
~ (6.5)

which means that vectors A ~×B ~ and B

~ ×A~ are antiparallel to each other and
that the cross product is not commutative but anti-commutative.
The magnitude of the cross product is given by
~ × B|
C = |A ~ = AB sin ϕ (6.6)

The cross product of two parallel vectors (ϕ = 00 ) vanishes. In this case

~·B
A ~ = AB sin 00 = 0 (6.7)

(6.7) implies that the cross product of a vector with itself vanishes.

6.3 Vector algebra - component form

Figure 6.3: (a) Unit vectors (b) Components of a vector.

6.3.1 Unit vectors

A unit vector is a vector of magnitude or length 1. In Cartesian coordinate
system, the the unit vectors along the x, y, z directions are represented by î, ĵ, k̂,

3
 6.3 Vector algebra - component form

respectively (see Figure 6.3(a)).

Since î, ĵ, k̂ are mutually perpendicular, from (6.2) we have

î · ĵ = ĵ · k̂ = k̂ · î = 0 (6.8)

And from (6.3), we have

î · î = ĵ · ĵ = k̂ · k̂ = 1 (6.9)

Any vector A ~ can be expressed in terms of the components Ax , Ay , Az (see

Figure 6.3(b)):

~ = Ax î + Ay ĵ + Az k̂
A (6.10)

where
~ · î, Ay = A
Ax = A ~ · ĵ, Az = A
~ · k̂ (6.11)

6.3.2 Vector operations

1. To add vectors, add like components:
~+B
A ~ = (Ax + Bx )î + (Ay + By )ĵ + (Az + Bz )k̂ (6.12)

2. To multiply by a scalar, multiply each component with the scalar:

~ = aAx î + aAy ĵ + aAz k̂
aA (6.13)

3. To calculate the dot product, mulitply like components, and then add:
~·B
A ~ = Ax Bx + Ay By + Az Bz (6.14)

This also follows from the properties of the unit vectors given by (6.8)
and 6.9.
In particular,
~·A
A ~ = A2x + A2y + A2z = A2 (6.15)

4. To calculate the cross product, form the determinant whose first row is

4
 6.4 Special vectors

~ and third row is B:

î, ĵ, k̂, second row is A, ~

î ĵ k̂
~ ~
A × B = Ax Ay Az
Bx By Bz
= (Ay Bz − Az By )î + (Az Bx − Ax Bz )ĵ + (Ax By − Ay Bx )k̂ (6.16)

6.4 Special vectors

1. Position vector (as shown in Figure 6.4) gives the position (x, y, z) of a
particle at any instant of time:

~r = xî + y ĵ + z k̂ (6.17)

Figure 6.4: Position vector.

The position vector has the magnitude

p
r = x2 + y 2 + z 2 (6.18)

and points in the direction given by the unit vector

~r
r̂ = (6.19)
r

5
 6.4 Special vectors

2. Infinitesimal displacement vector (d~l or d~r) represents a small dis-

placement from the position (x, y, z) to the position (x+dx, y+dy, z+dz):

d~l (or d~r) = dxî + dy ĵ + dz k̂ (6.20)

3. Del or nabla (∇) operator is the equivalent of the derivative in three

dimensions.
In one dimension, the derivative of a function f = f (x) is given by
 
df
df = dx (6.21)
dx

In three dimensions, for a function f = f (x, y, z),

∂f ∂f ∂f
df = dx + dy + dz (6.22)
∂x ∂y ∂z

which can be rearranged as

 
∂f ∂f ∂f 
= î + ĵ + k̂ dxî + dy ĵ + dz k̂
∂x ∂y ∂z
= ∇f · d~l (6.23)

where
∂f ∂f ∂f
∇f = î + ĵ + k̂ (6.24)
∂x ∂y ∂z
is a vector quantity called the gradient of f , with
∂ ∂ ∂
∇= î + ĵ + k̂ (6.25)
∂x ∂y ∂z

Geometrical interpretation of the gradient:

From (6.23),

df = ∇f · d~l = |∇f ||d~l| cos ϕ (6.26)

is maximum for ϕ = 0, i.e., when ∇f and d~l are in the same direction.
Thus, the gradient points in the direction of maximum increase of the function
f . The magnitude |∇f | gives the slope (rate of increase) along this direction.

6
 6.5 Operations of the ∇ operator

If ∇f = 0 at (x, y, z), then df = 0 for all small displacements about the

point (x, y, z). Then it is a stationary point (maximum, minimum, saddle point
or shoulder) of the function f .
Gradient of the position vector ~r:
∂r ∂r ∂r
∇r = î + ĵ + k̂
∂x ∂y ∂z
~r
= = r̂ (6.27)
r

6.5 Operations of the ∇ operator

The ∇ operator can operate in the three following ways:

1. on a scalar function f =⇒ ∇f (the gradient, discussed above in Section

6.4 and given by (6.24)).

2. on a vector function F~ , via the dot product =⇒ ∇ · F~ (the divergence).

3. on a vector function F~ , via the cross product =⇒ ∇ × F~ (the curl).

6.5.1 Divergence
  
~ ∂ ∂ ∂ 
∇·F = î + ĵ + k̂ · Fx î + Fy ĵ + Fz k̂
∂x ∂y ∂z
∂Fx ∂Fy ∂Fz
= + + (6.28)
∂x ∂y ∂z
which is a scalar quantity.

Geometrical interpretation of the divergence: It is a measure of how

much the vector F~ spreads out (i.e., diverges) from the point in question.
Examples of highly diverging functions are: flow of water from a tap (positive
divergence), and into a drain sink (negative divergence).
If ∇ · F~ = 0, then F~ is a solenoidal function.

7
 6.5 Operations of the ∇ operator

6.5.2 Curl

î ĵ k̂

∂ ∂ ∂
∇ × F~ =
∂x ∂y ∂z

F x Fy Fz

     
∂Fz ∂Fy ∂Fx ∂Fz ∂Fy ∂Fx
= − î + − ĵ + − k̂ (6.29)
∂y ∂z ∂z ∂x ∂x ∂y

Geometrical interpretation of the divergence: It is a measure of how

much the vector F~ swirls (curls) around the point in question.
Example of a function with a large curl is the motion of an object around a
circle.
If ∇ × F~ = 0, then F~ is a irrotational function.

Example 1: Find the divergence and curl of the position vector.

Solution: Position vector, ~r = xî + y ĵ + z k̂.

Divergence of ~r,
∂x ∂y ∂z
∇ · ~r = + + =1+1+1=3
∂x ∂y ∂z
Hence, ~r has a positive divergence, as expected.

Curl of ~r,

î ĵ k̂

∂ ∂ ∂
∇ × ~r = =0
∂x ∂y ∂z

x y z
Hence, ~r is irrotational.

~ = −y î + xĵ.
Example 2: Find the divergence and curl of the function A

8
 6.5 Operations of the ∇ operator

Solution:
Divergence of A,~
∇·A ~ = ∂(−y) + ∂x + ∂(0) = 0
∂x ∂y ∂z
Hence, ~r is solenoidal.

~
Curl of A,

î ĵ k̂

~ = ∂ ∂ ∂ = 0î + 0ĵ + (1 − (−1))k̂ = 2k̂

∇×A
∂x ∂y ∂z

−y x 0

Tip: Try to plot and check what the function above looks like, and if the
obtained values of divergence and curl make sense.

6.5.3 Second derivatives

1. Divergence of gradient

∂ 2f ∂ 2f ∂ 2f
∇ · (∇f ) = 2
+ 2 + 2 = ∇2 f
∂x ∂y ∂z

∇2 is called the Laplacian. Laplacian of a scalar f is also a scalar.

Laplacian of a vector F~ , ∇2 F~ = ∇2 Fx î + ∇2 Fy ĵ + ∇2 Fz k̂.

2. Curl of gradient ∇ × (∇f ) is always zero.

3. Gradient of divergence ∇(∇ · F~ ) 6= ∇2 F~ .

4. Divergence of curl ∇ · (∇ × F~ ) is always zero.

5. Curl of curl ∇ × (∇ × F~ ) = ∇(∇ · F~ ) − ∇2 F~ .

Tip: Check that the above statements are true.

9
 6.6 Integral calculus

6.6 Integral calculus

6.6.1 Line (or path) integral
Line (or path) integral of a vector F~ is defined as
Z ~b
F~ · d~l (6.30)
~a

where d~l is the infinitesimal displacement vector along a path from the point ~a
to the point ~b.
For a closed loop, i.e., ~b = ~a, the line integral is represented by
I ~b
F~ · d~l (6.31)
~a

For example, work done by a force F~ in moving an object from the point the
R ~b
point a to the point b is given by W = ~a F~ · d~l.

6.6.2 Surface integral

Surface integral of a vector F~ over a specified surface S is defined as
Z
F~ · d~a (6.32)
S

where d~a is an infinitesimal patch of area with direction perpendicular to the

surface.
In Cartesian coordinates,

d~a = dxdy k̂ or
= dydz î or
= dzdxĵ (6.33)

For a closed surface (e.g., a balloon, a hollow sphere), the surface integral is
represented by
I
F~ · d~a (6.34)
S

If F~ represents the flow of a quantity, then F~ · d~a represents the total

R
S

10
 6.6 Integral calculus

amount per unit time passing through the surface, and is called flux of the
vector F~ .

6.6.3 Volume integral

Volume integral of a scalar f is defined as
Z
f dV (6.35)
V

where dV is an infinitesimal volume element within the volume V .

In Cartesian coordinates,

dV = dxdydz (6.36)

For a vector F~ ,
Z Z Z Z
F~ dV = Fx dV î + Fy dV ĵ + Fz dV k̂ (6.37)
V V V V

6.6.4 Fundamental theorem for gradients

In one dimension, for a scalar function f (x),
Z b 
df
dx = f (b) − f (a) (6.38)
a dx

Similarly, in three dimensions, for a scalar function f (x, y, z),

Z ~b
∇f · d~l = f (~b) − f (~a) (6.39)
~a

Corollaries
R ~b
1. The integral ~a ∇f · d~l is path-independent (as it is dependent only on
the values of the function at the endpoints of the path, as evident from
(6.39)).

2. For a closed loop,

I
∇f · d~l = 0 (6.40)

11
 6.6 Integral calculus

6.6.5 Fundamental theorem for divergences

Integral of the divergence of a function over a volume V is equal to the value
of the function at the surface S enclosing the volume (or flux of the function).
Z I
(∇ · F~ )dV = F~ · d~a (6.41)
V S

This is known as Gauss’ theorem, or Green’s theorm, or divergence theorem.

6.6.6 Fundamental theorem for curls

Surface integral of the curl of a function is equal to the line integral of the
function around the periphery of the surface.
Z I
(∇ × F~ ) · d~a = F~ · d~l (6.42)
S

This is known as Stokes’ theorem.

********************* The End (of this Unit) *********************

12
 Electromagnetic theory
Part I: Maxwell’s equations
PH-1007 (Physics)
Dr. Gorky Shaw

8.1 Fundamental laws of electromagnetism

8.1.1 Gauss’s law in electrostatics
q
Electric flux through any surface enclosing a charge q is .
0
I
~ · d~a = q
E (8.1)
0
S

Applying Gauss’s theorem,

I Z
~ · d~a =
E ~
(∇ · E)dV (8.2)
S V

If ρ is the charge density, then

Z
q= ρdV (8.3)
V

From (8.1), (8.2), and (8.3),

Z Z
~ ρ
(∇ · E)dV = dV
0
V V
~ = ρ
=⇒ ∇ · E (8.4)
0

(8.1) and (8.4) are Gauss’s law in electrostatics in integral form and differen-
tial form, respectively.

8.1.2 Gauss’s law in magnetostatics

I
~ · d~a = 0
B (8.5)
S
 8.1 Fundamental laws of electromagnetism
or
~ =0
∇·B (8.6)

(8.5) and (8.6) are Gauss’s law in magnetostatics in integral form and differ-
ential form, respectively.
Unlike (8.1) and (8.4), the RHS of (8.5) and (8.6) are zero. This indicates,
unlike electric charges, magnetic charges or magnetic monopoles do not exist.

8.1.3 Faraday’s law of electromagnetic induction

Changing magnetic flux induces an emf in a closed circuit, which opposes the
change in magnetic flux.
dΦ
ε=− (8.7)
dt
Now,
Z
Φ= ~ · d~a
B (8.8)

Therefore, from (8.7) and (8.8),

d
Z Z ~
∂B
ε=− ~ · d~a = −
B · d~a (8.9)
dt ∂t
(In the last step of (8.9), the total derivative is replaced by the partial derivative
because in general B ~ is a function of several variables, and not just t).
Also,
I Z
ε= E ~ · d~l = (∇ × E) ~ · d~a (8.10)

Therefore, from (8.9) and (8.10),

I Z Z ~
∂B
~ · d~l =
E ~ · d~a = −
(∇ × E) · d~a (8.11)
∂t

2
 8.1 Fundamental laws of electromagnetism

Since this applies to any surface, we have

~
∂B
~
∇×E =− (8.12)
∂t

(8.11) and (8.12) are Faraday’s law of electromagnetic induction in integral

form and differential form, respectively.

8.1.4 Ampere’s law

Line integral of magnetic flux around a closed loop is equal to µ0 times the
current enclosed by the loop.
I
~ · d~l = µ0 I
B (8.13)

If J~ is the current density, then

Z
I= J~ · d~a (8.14)

with the integral taken over any surface bounded by the loop. And
I Z
B~ · d~l = (∇ × B)
~ · d~a (8.15)

Therefore, from (8.13), (8.14) and (8.15),

Z Z
(∇ × B)~ · d~a = µ0 J~ · d~a

~ = µ0 J~
=⇒ ∇ × B (8.16)

(8.13) and (8.16) are Ampere’s law in integral form and differential form,
respectively.

3
 8.1 Fundamental laws of electromagnetism

8.1.5 Equation of continuity

If charge q is enclosed in a volume V by a surface S, then current flowing out
of the surface is given by
∂q
I=− (8.17)
∂t
Now, from (8.14), and applying Gauss’s divergence theorem,
I Z
I = J~ · d~a = (∇ · J)dV
~ (8.18)
S V

and from (8.3),

Z
q= ρdV (8.19)
V

Therefore, from (8.17), (8.18) and (8.19),

Z Z
~ ∂
(∇ · J)dV =− ρdV
∂t
V
Z Z V
~ ∂ρ
=⇒ (∇ · J)dV =− dV
∂t
V V
Z  
∂ρ
=⇒ ∇ · J~ + dV = 0 (8.20)
∂t
V

Since this applies to any volume, we have

∂ρ
∇ · J~ + =0 (8.21)
∂t

(8.21) is known as the continuity equation, equation of continuity, or charge

conservation law.

8.1.6 Maxwell’s modification of Ampere’s law

The problem:

4
 8.1 Fundamental laws of electromagnetism

Taking divergence of Ampere’s law given by (8.15), we get

~ = ∇ · (µ0 J)
∇ · (∇ × B) ~
~ = µ0 (∇ · J)
or, ∇ · (∇ × B) ~ (8.22)

We know that divergence of curl of any vector is always zero. Therefore,

~ =0
∇ · (∇ × B) (8.23)

From (8.22) and (8.23), it follows that

∇ · J~ = 0 (8.24)

But from the continuity equation (8.21), it follows that

∂ρ
∇ · J~ = − 6= 0 for non-steady currents. (8.25)
∂t

Ampere’s law in the form of ∇ · J~ = 0 given by (8.24) is valid only when the
charge density ρ is static and does not change with time.

The solution by Maxwell:

From Gauss’s law of electrostatics (8.4) and the continuity equation (8.21),
we get
!
∂ρ ∂ ~
∇ · J~ = − = − (0 ∇ · E)~ = −∇ · 0 ∂ E
∂t ∂t ∂t
!
∂E~
=⇒ ∇ · J~ + ∇ · 0 =0
∂t
!
∂E~
=⇒ ∇ · J~ + 0 =0 (8.26)
∂t

Replacing the RHS of (8.22) with LHS of (8.26), we get

!
~
~ = µ0 ∇ · J~ + 0 ∂ E = 0
∇ · (∇ × B) (8.27)
∂t

which is consistent.

5
 8.2 Maxwell’s equations of electromagnetism

Thus, Ampere’s law is modified as

~
∂E
~ ~
∇ × B = µ0 J + µ0 0 (8.28)
∂t
Taking integral of (8.28) over a surface bounded by the current loop in ques-
tion, we get
!
Z Z ~
∂E
(∇ × B)~ · d~a = µ0 J~ + µ0 0 · d~a
∂t
S S
!
I Z ~
∂E
or, ~ · d~l =
B µ0 J~ + µ0 0 · d~a (8.29)
∂t
S

~ does not change with time, we get back the original Ampere’s law
When E
~ = µ0 J~
∇×B (8.30)

(8.28) and (8.29) are modified Ampere’s law in differential form and integral
form, respectively.

8.1.7 Displacement current

The modified Ampere’s law implies that not only a conduction current, but also
a changing electric field produces a magnetic field. The equivalent current due
to a changing electric field is known as displacement current. For example,
charging or discharging a capacitor produces a time-varying electric field, and
hence sets up a displacement current between the capacitor plates.
Conduction current involves flow of electrons, while displacement current
involves displacement of electrons in a time-varying electric field.

8.2 Maxwell’s equations of electromagnetism

8.2.1 Maxwell’s equations in differential form
The fundamental laws of electromagnetism together are known as Maxwell’s
equations of electromagnetism. The four Maxwell’s equations in differential

6
 8.2 Maxwell’s equations of electromagnetism

form are:
~ = ρ
∇·E (Gauss’s law of electrostatics (8.4)) (8.31)
0
~ = 0 (Gauss’s law of magnetostatics (8.6))
∇·B (8.32)
∂B~
~
∇×E =− (Faraday’s law of electromagnetic induction (8.12)) (8.33)
∂t
~
~ = µ0 J~ + µ0 0 ∂ E (Modified Ampere’s law (8.28))
∇×B (8.34)
∂t

In free space, i.e., in a region where ρ = 0 and J~ = 0, the Maxwell’s

equations are:
~ =0
∇·E (8.35)
~ =0
∇·B (8.36)
~
∇×E ~ = − ∂B (8.37)
∂t
~
∂E
~
∇ × B = µ0 0 (8.38)
∂t

8.2.2 Maxwell’s equations in integral form

The four Maxwell’s equations in integral form are:
I
~ · d~a = q
E (8.39)
I 0
~ · d~a = 0
B (8.40)
I Z ~
E~ · d~l = − ∂ B · d~a (8.41)
∂t !
I Z ~
∂E
B~ · d~l = µ0 J~ + µ0 0 · d~a (8.42)
∂t

In free space, i.e., in a region where ρ = 0 and J~ = 0, the Maxwell’s equations

7
 8.2 Maxwell’s equations of electromagnetism

are:
I
~ · d~a = 0
E (8.43)
I
~ · d~a = 0
B (8.44)
!
I Z ~
∂B
~ · d~l = −
E · d~a (8.45)
∂t
!
I Z ~
∂E
~ · d~l =
B µ0 0 · d~a (8.46)
∂t

~ and H
8.2.3 Maxwell’s equations in terms of D ~
~ is the electric displacement vector D
Associated with the electric field vector E ~
given by
~ = 0 E
D ~ (8.47)

~ is the magnetic field strength

And, associated with the magnetic field vector B
~ given by
vector H
~
~ = B
H (8.48)
µ0
Using these vectors, the Maxwell’s equations can be written in a more sym-
metric form.
Differential form:
~ =ρ
∇·D (8.49)
~ =0
∇·B (8.50)
~
~ = − ∂B
∇×E (8.51)
∂t
~
~ = J~ + ∂ D
∇×H (8.52)
∂t

8
 8.2 Maxwell’s equations of electromagnetism

Differential form in free space:

~ =0
∇·D (8.53)
~ =0
∇·B (8.54)
~
~ = − ∂B
∇×E (8.55)
∂t
~
~ = ∂D
∇×H (8.56)
∂t
Integral form:
I
~ · d~a = q
D (8.57)
I
~ · d~a = 0
B (8.58)
I Z
∂B~
~ · d~l = −
E · d~a (8.59)
∂t !
I Z
∂ D~
~ · d~l =
H J~ + · d~a (8.60)
∂t

Integral form in free space:

I
D~ · d~a = 0 (8.61)
I
B~ · d~a = 0 (8.62)
I Z ~
~ · d~l = − ∂ B · d~a
E (8.63)
∂t !
I Z
∂D~
H~ · d~l = · d~a (8.64)
∂t

********************* The End (of this Unit) *********************

9
 Electromagnetic theory
Part II: Electromagnetic waves
PH-1007 (Physics)
Dr. Gorky Shaw

9.1 Poynting theorem: electromagnetic energy

Maxwell’s third equation is:
~
~ = − ∂B
∇×E (9.1)
∂t
~ we get
Taking dot product of (9.1) with B,

~
~ · (∇ × E)
B ~ · ∂B
~ = −B (9.2)
∂t
Maxwell’s fourth equation is:
~
~ = µ0 J~ + µ0 0 ∂ E
∇×B (9.3)
∂t
~ we get
Taking dot product of (9.3) with E,

~
~ · (∇ × B)
E ~ = µ0 (E
~ · J) ~ · ∂E
~ + µ0 0 E (9.4)
∂t
Using (9.2) and (9.4), and the vector identity

~ × B)
∇ · (E ~ =B
~ · (∇ × E)
~ −E
~ · (∇ × B)
~ (9.5)

we get
! !
~
∂B ~
~ × B)
∇ · (E ~ = ~
−B · − µ0 (E ~ · J) ~ · ∂E
~ + µ0 0 E
∂t ∂t
!
~ ~
~ × B)
∇ · (E ~ · ∂ B + µ0 0 E
~ =− B ~ · ∂ E − µ0 (E ~ · J)
~ (9.6)
∂t ∂t
 9.1 Poynting theorem: electromagnetic energy
Now,
~
~ · ∂ B = Bx ∂Bx + By ∂By + Bz ∂Bz
B
∂t ∂t ∂t ∂t
1∂ 2
= (Bx + By2 + Bz2 )
2 ∂t
1∂ 2
= (B ) (9.7)
2 ∂t
Similarly,
~
~ · ∂ E = 1 ∂ (E 2 )
E (9.8)
∂t 2 ∂t
Using (9.7) and (9.8) in (9.6), we get

∇ · (E ~ = − 1 ∂ (B 2 + µ0 0 E 2 ) − µ0 (E
~ × B) ~ · J)
~ (9.9)
2 ∂t
Rearranging,
2
 
~ · J~ = − 1 ∂ 2 B 1 ~ × B)
~
E 0 E + − ∇ · (E (9.10)
2 ∂t µ0 µ0

Integrating over a volume V bound by a surface S, and using the divergence

theorem, we get
Z  2
Z  Z
~ · J)dV
~ ∂ 1 2 B 1 ~ × B))dV
~
(E =− 0 E + dV − (∇ · (E
∂t 2 µ0 µ0
V V V
Z  2
Z  I
~ · J)dV
~ ∂ 1 2 B 1 ~ × B)
~ · d~a (9.11)
or, (E =− 0 E + dV − (E
∂t 2 µ0 µ0
V V S

(9.11) is known as the Poynting theorem.

Significance of the terms in (9.11):
~ · J)dV
~
R
1. The term (E represents the total power (work done per unit time,
V
dW
) in a volume V .
dt
2. The first term on the RHS is the rate of change of total energy stored
(sum of energies stored in the electric and magnetic fields) in volume V .

3. The second term on the RHS represents the rate at which energy is trans-

2
 9.2 Electromagnetic waves in vacuum

ported out of the volume V , across its boundary surface, by the electro-
magnetic fields.

The energy per unit time, per unit area, transported by the fields is called
~
the Poynting vector, S:

~ = 1 (E
S ~ × B)
~ (9.12)
µ0

9.2 Electromagnetic waves in vacuum

~ and B
9.2.1 Wave equations for the E ~ fields in vacuum

Vacuum can be treated as a free space with no charges and currents, i.e., ρ = 0
and J~ = 0.
We have Maxwell’s third equation:
~
~ = − ∂B
∇×E (9.13)
∂t
Applying the curl to the LHS of (9.13), we get

~ = ∇(∇ · E)
∇ × (∇ × E) ~ − ∇2 E
~
~
= 0 − ∇2 E
~
= −∇2 E (9.14)

~ = 0 in free space (Maxwell’s first equation).

since ∇ · E
Applying the curl to the RHS of (9.13), we get
!
~
∂B ∂
∇× − = − (∇ × B) ~
∂t ∂t
!
∂ ~
∂E
=− µ0 0
∂t ∂t
~
∂ 2E
= −µ0 0 2 (9.15)
∂t
where we have used Maxwell’s fourth equation in free space (9.3).

3
 9.2 Electromagnetic waves in vacuum

Therefore, from (9.13), (9.14), and (9.15), we get

~
∂ 2E
2~
∇ E = µ0 0 2 (9.16)
∂t
Similarly, it can be shown that
~
∂ 2B
2~
∇ B = µ0 0 2 (9.17)
∂t
These are wave equations for the electric field and magnetic field components,
respectively, of electromagnetic waves (EM waves) in free space.

~ and B
In general, each Cartesian component of the E ~ fields satisfies the
three-dimensional wave equation

1 ∂ 2y
∇2 y = (9.18)
v 2 ∂t2
where velocity of the wave is given by
1
v=√ (9.19)
µ0 0

Now,
µ0
= 10−7 N/A2 and
4π
1
= 9 × 109 Nm2 /C2 (9.20)
4π0
Therefore,

10−7 1
µ0 0 = = (9.21)
9 × 109 9 × 1016

Hence,
1
v=√ = 3 × 108 m/s = c (the speed of light). (9.22)
µ0 0

That is, electromagnetic waves propagate in free space with the speed of light.

4
 9.2 Electromagnetic waves in vacuum

Thus, we can rewrite the wave equations (9.16) and (9.17) as

~
1 ∂ 2E
2~
∇E= 2 2 (9.23)
c ∂t
and
~
1 ∂ 2B
~ =
∇2 B (9.24)
c2 ∂t2

9.2.2 Transverse of nature of electromagnetic waves in vacuum

The wave equations given by (9.23) and (9.24) have the general solutions

~ r, t) = E~0 ei(~k·~r−ωt)
E(~
~ r, t) = B
and B(~ ~0 ei(~k·~r−ωt) (9.25)

where E~0 and B~0 are the complex amplitudes. The physical fields are their
real parts. And, ω = ck, where ~k is the propagation vector:

~k = kn̂ = 2π n̂
λ
2πν ω
= n̂ = n̂ (9.26)
c c
where n̂ is the unit vector along the direction of propagation of the EM waves.
Applying the divergence to the E ~ field given by (9.25), we get
 
~ ~ i(~k·~r−ωt)
∇ · E = ∇ · E0 e
 
∂ ∂ ∂
= î + ĵ + k̂ · (E0x î + E0y ĵ + E0z k̂)ei(kx x+ky y+kz z−ωt)
∂x ∂y ∂z
∂ 
i(kx x+ky y+kz z−ωt)
 ∂  i(kx x+ky y+kz z−ωt)

= E0x e + E0y e
∂x ∂y
∂  i(kx x+ky y+kz z−ωt)

+ E0z e (9.27)
∂z

5
 9.2 Electromagnetic waves in vacuum

The first term of RHS of (9.27)

= E0x (ikx )ei(kx x+ky y+kz z−ωt)

~
= ikx E0x e(k·~r−ωt) (9.28)

Therefore, RHS of (9.27)

~
= i[kx E0x + ky E0y + kz E0z ]e(k·~r−ωt)
~
= i~k · E~0 e(k·~r−ωt)
= i(~k · E)
~ (9.29)

Therefore, from (9.27) and (9.29),

~ = i(~k · E)
∇·E ~ =0
i.e., ~k · E
~ =0 (9.30)

~ = 0 in free space.
since ∇ · E
Similarly, it can be shown that
~k · B
~ =0 (9.31)

~ and the magnetic field

(9.30) and (9.31) imply that both the electric field E
B~ are perpendicular to the direction of propagation (which is along ~k). That
is, electromagnetic waves in free space are transverse in nature.

Similarly, it can also be shown that

~ = i(~k × E)
∇×E ~ (9.32)

Now, Maxwell’s third equation is

~
~ = − ∂B
∇×E (9.33)
∂t
and from (9.25),

~
∂B
− ~0 ei(~k·~r−ωt)
= iω B
∂t
~
= iω B (9.34)

6
 9.2 Electromagnetic waves in vacuum

Therefore, from (9.32), (9.33) and (9.34),

i(~k × E)
~ = iω B
~
=⇒ ~k × E~ = ωB~ (9.35)

Similarly, using Maxwell’s fourth equation in free space, i.e.,

~
~ = µ0 0 ∂ E
∇×B (9.36)
∂t
it can be shown that
~k × B
~ = ωµ0 0 E
~ (9.37)

~ and B
(9.35) and (9.37) imply that E ~ are also mutually perpendicular to
each other.

Taking only magnitude of (9.35), we get

kE = ωB
E ω
=⇒ = =c (9.38)
B k
Using B = µ0 H, we get
r
E µ0
Z0 = = µ0 c = = 376.72 Ω (9.39)
H 0

which is called the wave impedance of free space.

9.2.3 Energy in a plane electromagnetic wave

The energy per unit time, per unit area, transported by the electric and mag-
~
netic fields is given by the Poynting vector S:

~ = 1 (E
S ~ × B)
~
µ0
1 ~
= [E × (~k × E)]
~ (9.40)
µ0 ω

7
 9.2 Electromagnetic waves in vacuum

Using the vector identity

~ × (B
A ~ × C)
~ = (A
~ · C)
~ B~ − (A
~ · B)
~ C~ (9.41)

we get

~ = 1 [(E
S ~ ~k − (E
~ · E) ~ · ~k)E]
~
µ0 ω
E2 ~ ~ · ~k = 0)
= k (since E
µ0 ω
E2
= n̂ (using (9.26)) (9.42)
µ0 c
~
Over a complete cycle of an EM wave, average value of S:

~ = hE 2 i
hSi n̂ (9.43)
µ0 c
For a sinusoidal wave function given by

~ r, t) = E~0 ei(~k·~r−ωt)
E(~ (9.44)

we have

2 E02 2
hE i = = Erms (9.45)
2
where
E0
Erms = √ (9.46)
2
is the RMS value of the electric field.

Therefore,
2
~ Erms
hSi = n̂ (9.47)
µ0 c

8
 9.3 Wave equation in a charge-free non-conducting medium

9.2.4 Energy density

Energy stored per unit volume in an EM field is given by the sum of the energies
stored in the electric field and the magnetic field:

B2
 
1 2
U= 0 E + (9.48)
2 µ0

From (9.38), in free space,

E √
B= = µ0 0 E (9.49)
c
Using (9.49) in (9.48), we get
1
0 E 2 + 0 E 2 = 0 E 2


U= (9.50)
2
Therefore, average density per unit cycle of an EM wave,

hU i = 0 hE 2 i = 0 Erms
2
(9.51)

Therefore, from (9.47) and (9.51), we get

 
hSi 1 1
= n̂ = cn̂ since µ0 0 = 2 (9.52)
hU i µ0 0 c c

Hence,

hSi = chU in̂ (9.53)

or, energy flux = velocity of light × energy stored in the fields.

9.3 Wave equation in a charge-free non-conducting medium

In a charge-free, non-conducting medium, we still have ρ = 0 and J~ = 0.
But unlike vacuum or free space, we need to consider the permittivity () and
permeability (µ) of the medium, instead of those of free space (0 and µ0 ).
In this case, the wave equations for the electric and magnetic fields given by
(9.16) and (9.17) are modified to

~
∂ 2E
~ = µ
∇2 E (9.54)
∂t2

9
 9.4 Wave equations in terms of scalar and vector potentials

and
~
∂ 2B
2~
∇ B = µ 2 (9.55)
∂t
These are, however, valid only if  and µ are independent of position and time.
Speed of EM wave in the medium,
1
v=√ (9.56)
µ

Since we always have µ > µ0 and  > 0 , (9.56) indicates that we always have
v < c.

9.4 Wave equations in terms of scalar and vector potentials

9.4.1 Vector and scalar potentials

~ the vector potential F~ is defined by the condition

For a given vector field G,
~ = ∇ × F~
G (9.57)

We know that any vector function F~ satisfies the condition

∇ · (∇ × F~ ) = 0 (9.58)

~ admits a vector potential, then from (9.57) and (9.58),

Thus, if a vector G
we have
~ =0
∇·G (9.59)

~ = 0, then corresponding to G,
Conversely, if ∇ · G ~ we can always define a vector
potential F~ such that G~ = ∇ × F~ .

~ the scalar potential P is defined by the condition

For a given vector field G,
~ = −∇P
G (9.60)

10
 9.4 Wave equations in terms of scalar and vector potentials

9.4.2 Magnetic vector potential

~ = 0, we can define the magnetic vector potential A,
Since ∇ · B ~ such that

~ =∇×A
B ~ (9.61)

However, A~ is not unique.

~ we have
For a constant vector C,
~ + C)
∇ × (A ~ = (∇ × A)
~ + (∇ × C)
~ =B
~ +0=B
~ (9.62)

For a gradient of a scalar ∇f , we have

~ + ∇f ) = (∇ × A)
∇ × (A ~ + (∇ × (∇f )) = B
~ +0=B
~ (9.63)

~ (A
(9.61), (9.62), and (9.63) imply that A, ~ + C),
~ and (A
~ + ∇f ) all give the
~
same B.

So, the magnetic vector potential is not unique, but arbitrary to the extent
of addition of a constant vector or the gradient of a scalar.
This arbitrariness allows one to choose a convenient value of A~ for mathe-
matical simplicity.

9.4.3 The scalar potential

Maxwell’s third equation gives
~
~ + ∂B = 0
(∇ × E)
∂t
∂
~ + (∇ × A) ~ =0
=⇒ (∇ × E)
∂t !
~
=⇒ ∇ × E ~ + ∂A = 0 (9.64)
∂t

Since ∇ × (∇f ) for any scalar function f , we can define a scalar function φ
such that
!
~
∇× E ~ + ∂ A = −∇ × (∇φ) = 0
∂t
~
∂A
~
or, E + = −∇φ (9.65)
∂t

11
 9.4 Wave equations in terms of scalar and vector potentials

~ is defined
φ is called the scalar potential, and from (9.65), the electric field E
in terms of the scalar potential (and the vector potential) as

~
~ = −∇φ − ∂ A
E (9.66)
∂t
~
∂A
For a time-independent field, = 0, so we get
∂t
~ = −∇φ
E (9.67)

In this case, φ is the electrostatic potential.

9.4.4 The Lorenz gauge condition

The scalar and vector potentials are not unique. Any value of A ~ and φ are
~ and B
acceptable as long as the correct E ~ fields are obtained from them by the
relations (9.61) and (9.66).
If we make the transformations

A ~ − ∇f and φ → φ + ∂f
~→A (9.68)
∂t

where f (~r, t) is any scalar function, then values of E~ and B

~ remain unchanged
or invariant.
Check: The above transformations transform the E ~ and B~ fields as
 
~ ∂f ∂ ~ 
E → −∇ φ + − A − ∇f
∂t ∂t
∂A~
= −∇φ − (the other two terms cancel out)
∂t
=E ~ (9.69)

and

~ → ∇ × (A
B ~ − ∇f )
=∇×A ~+0
~
=B (9.70)

12
 9.4 Wave equations in terms of scalar and vector potentials

~ and B
That is, the E ~ fields remain unchanged under the transformations of φ
~ given by (9.68)
and A

Transformations of the potentials that leave the fields invariant are called
Gauge transformations.
Various gauge conditions can be applied on the potentials depending on the
requirement and convenience.
Two important and frequently used Gauge conditions are: the Lorenz gauge
condition given by

~+ 1 ∂φ
∇·A =0 (9.71)
c2 ∂t
or

~ + µ0 0 ∂φ = 0
∇·A (9.72)
∂t
and the Coulomb gauge condition given by
~=0
∇·A (9.73)

~
9.4.5 Wave equations in terms of φ and A
In free space, we have
~ =0
∇·E (9.74)

~ in terms of φ and A,
Expressing E ~ we get
!
~
∂A
∇· −∇φ − =0
∂t
∂ ~ =0
=⇒ ∇2 φ + (∇ · A) (9.75)
∂t
Using the Lorenz gauge condition (9.71) in (9.75), we get
 
∂ 1 ∂φ
∇2 φ + − 2 =0
∂t c ∂t
2 1 ∂ 2φ
=⇒ ∇ φ − 2 2 = 0 (9.76)
c ∂t

13
 9.4 Wave equations in terms of scalar and vector potentials

(9.76) gives

1 ∂ 2φ
∇2 φ =
c2 ∂t2
2 ∂ 2φ
or, ∇ φ = µ0 0 2 (9.77)
∂t
which is the wave equation in terms of the scalar potential.

Maxwell’s fourth equation gives

~
∇×B ~ = µ0 0 ∂ E
∂t !
~
~ = µ0 0 ∂ −∇φ − ∂ A
=⇒ ∇ × (∇ × A)
∂t ∂t
2~
=⇒ ∇(∇ · A) ~ = −µ0 0 ∂ (∇φ) − µ0 0 ∂ A
~ − ∇2 A
∂t ∂t2

∂φ

∂ 2A~
~
∇ ∇ · A + µ0 0 2~
− ∇ A = −µ0 0 2 (9.78)
∂t ∂t

From the Lorenz condition (9.72), the first term of (9.78) is zero. Therefore,
we get

~
∂ 2A
~ = µ0 0
∇2 A
∂t2
1 ∂ 2A~
2~
or, ∇ A = 2 2 (9.79)
c ∂t
which is the wave equation in terms of the magnetic vector potential.

~ have the same form as that

Thus, the wave equations in terms of φ and A
~ and B.
for E ~

********************* The End (of this Unit) *********************

14
 Electromagnetic theory
Part II: Electromagnetic waves
PH-1007 (Physics)
Dr. Gorky Shaw

9.1 Electromagnetic waves in vacuum

~ and B
9.1.1 Wave equations for the E ~ fields in vacuum

Vacuum can be treated as a free space with no charges and currents, i.e., ρ = 0
and J~ = 0.
We have Maxwell’s third equation:
~
~ = − ∂B
∇×E (9.1)
∂t
Applying the curl to the LHS of (9.1), we get

~ = ∇(∇ · E)
∇ × (∇ × E) ~ − ∇2 E
~
~
= 0 − ∇2 E
~
= −∇2 E (9.2)

~ = 0 in free space (Maxwell’s first equation).

since ∇ · E
Applying the curl to the RHS of (9.1), we get
!
~
∂B ∂
∇× − = − (∇ × B) ~
∂t ∂t
!
∂ ~
∂E
=− µ0 0
∂t ∂t
~
∂ 2E
= −µ0 0 (9.3)
∂t2
where we have used Maxwell’s fourth equation in free space (??).

Therefore, from (9.1), (9.2), and (9.3), we get

~
∂ 2E
~ = µ0 0
∇2 E (9.4)
∂t2
 9.1 Electromagnetic waves in vacuum

Similarly, it can be shown that

~
∂ 2B
2~
∇ B = µ0 0 2 (9.5)
∂t
These are wave equations for the electric field and magnetic field components,
respectively, of electromagnetic waves (EM waves) in free space.

~ and B
In general, each Cartesian component of the E ~ fields satisfies the
three-dimensional wave equation

21 ∂ 2y
∇y= 2 2 (9.6)
v ∂t
where velocity of the wave is given by
1
v=√ (9.7)
µ0 0

Now,
µ0
= 10−7 N/A2 and
4π
1
= 9 × 109 Nm2 /C2 (9.8)
4π0
Therefore,

10−7 1
µ0 0 = = (9.9)
9 × 109 9 × 1016

Hence,
1
v=√ = 3 × 108 m/s = c (the speed of light). (9.10)
µ0 0

That is, electromagnetic waves propagate in free space with the speed of light.
Thus, we can rewrite the wave equations (9.4) and (9.5) as

~
1 ∂ 2E
2~
∇E= 2 2 (9.11)
c ∂t

2
 9.1 Electromagnetic waves in vacuum

and
~
1 ∂ 2B
2~
∇B= 2 2 (9.12)
c ∂t

9.1.2 Transverse of nature of electromagnetic waves in vacuum

The wave equations given by (9.11) and (9.12) have the general solutions

~ r, t) = E~0 ei(~k·~r−ωt)
E(~
~ r, t) = B
and B(~ ~0 ei(~k·~r−ωt) (9.13)

where E~0 and B~0 are the complex amplitudes. The physical fields are their
real parts. And, ω = ck, where ~k is the propagation vector:

~k = kn̂ = 2π n̂
λ
2πν ω
= n̂ = n̂ (9.14)
c c
where n̂ is the unit vector along the direction of propagation of the EM waves.
Applying the divergence to the E ~ field given by (9.13), we get
 
~ ~ i(~k·~r−ωt)
∇ · E = ∇ · E0 e
 
∂ ∂ ∂
= î + ĵ + k̂ · (E0x î + E0y ĵ + E0z k̂)ei(kx x+ky y+kz z−ωt)
∂x ∂y ∂z
∂ 
i(kx x+ky y+kz z−ωt)
 ∂  i(kx x+ky y+kz z−ωt)

= E0x e + E0y e
∂x ∂y
∂  i(kx x+ky y+kz z−ωt)

+ E0z e (9.15)
∂z
The first term of RHS of (9.15)

= E0x (ikx )ei(kx x+ky y+kz z−ωt)

~
= ikx E0x e(k·~r−ωt) (9.16)

3
 9.1 Electromagnetic waves in vacuum

Therefore, RHS of (9.15)

~
= i[kx E0x + ky E0y + kz E0z ]e(k·~r−ωt)
~
= i~k · E~0 e(k·~r−ωt)
= i(~k · E)
~ (9.17)

Therefore, from (9.15) and (9.17),

~ = i(~k · E)
∇·E ~ =0
i.e., ~k · E
~ =0 (9.18)

~ = 0 in free space.
since ∇ · E
Similarly, it can be shown that
~k · B
~ =0 (9.19)

~ and the magnetic field

(9.18) and (9.19) imply that both the electric field E
B~ are perpendicular to the direction of propagation (which is along ~k). That
is, electromagnetic waves in free space are transverse in nature.

Similarly, it can also be shown that

~ = i(~k × E)
∇×E ~ (9.20)

Now, Maxwell’s third equation is

~
∂B
~
∇×E =− (9.21)
∂t
and from (9.13),

~
∂B
− ~0 ei(~k·~r−ωt)
= iω B
∂t
~
= iω B (9.22)

Therefore, from (9.20), (9.21) and (9.22),

i(~k × E)
~ = iω B
~
=⇒ ~k × E~ = ωB~ (9.23)

4
 9.1 Electromagnetic waves in vacuum

Similarly, using Maxwell’s fourth equation in free space, i.e.,

~
∂E
~
∇ × B = µ0 0 (9.24)
∂t
it can be shown that
~k × B
~ = ωµ0 0 E
~ (9.25)

~ and B
(9.23) and (9.25) imply that E ~ are also mutually perpendicular to
each other.

Taking only magnitude of (9.23), we get

kE = ωB
E ω
=⇒ = =c (9.26)
B k
Using B = µ0 H, we get
r
E µ0
Z0 = = µ0 c = = 376.72 Ω (9.27)
H 0

which is called the wave impedance of free space.

5
 Magnetic properties of matter
PH-1007 (Physics)
Dr. Gorky Shaw

10.1 Introduction
All magnetic phenomena are due to electric charges in motion. Motion of
electrons - orbital motion about the nucleus and spin motion about their own
axes, produce their own magnetic field.
Electrons possess an orbital magnetic moment as well as a spin magnetic
moment, which are randomly oriented in normal state. Hence, the resultant
magnetic moment of an atom is zero.
The moments tend to align in a certain way in response to an external mag-
netic field, and the material is said to be magnetized.

10.1.1 Parameters associated with magnetism

1. Magnetic field strength (H ): represents the strength of an externally
applied magnetic field, independent of the internal response of the ma-
terial / medium to which the field is applied. Magnetic lines of force
that diffuse through the medium are produced by the applied field. Unit:
A/m.

2. Magnetic flux density (B): the number of lines of force crossing per
unit area of the medium. Unit: Wb/m2 or Tesla.

3. Magnetization (M ): A magnetized material consists of many tiny mag-

netic dipoles, with a net alignment in some direction. The magnetic dipole
moment per unit volume is known as magnetization. It represents the in-
ternal response of a material to an external applied magnetic field. Unit:
A/m.
The B, H, and M fields are related as

B = µ0 (H + M ) (10.1)

where µ0 is the permeability of free space.

4. Magnetic susceptibility (χ): By definition,

M ∝H
or, M = χH (10.2)
 10.2 Classification of magnetic materials
where χ is called the magnetic susceptibility of the material. Materials
that obey (10.2) are called linear media.
5. Relative permeability (µr ): From (10.1),

B∝H
or, B = µH (10.3)

where µ is called the permeability of the material.

In vacuum (or in air), µ = µ0 , and the magnetic flux density B0 is given
by

B0 = µ0 H (10.4)

From (10.3) and (10.4), we have

B µ
= = µr (10.5)
B0 µ0
where µr is called the relative permeability of the material.

10.1.2 Relationship between µr and χ

From (10.1), (10.2) and (10.3), we get

B = µ0 (H + M )
=⇒ µH = µ0 (H + χH)
µ
=⇒ =1+χ
µ0
=⇒ µr =1+χ (10.6)

10.2 Classification of magnetic materials

Materials are classified according to the values of χ and µr into three broad
categories:
1. Diamagnetic: χ < 0 and µr < 1.
2. Paramagnetic: χ > 0 and µr > 1.
3. Ferromagnetic: χ >> 0 and µr >> 1.

2
 10.2 Classification of magnetic materials

10.2.1 Diamagnetic materials

In diamagnetic substances, the orbital and spin magnetic moments of atoms
are randomly oriented and add vectorially to zero.
The atoms acquire an induced dipole moment when placed in an external
magnetic field. The substance acquires a net magnetic moment in a direction
opposite to the magnetic field.
Diamagnetic materials lose their magnetization as soon as the applied mag-
netic field is removed.
They are feebly repelled by a magnet, and field lines tend to avoid the sub-
stance (as shown in Figure 10.1(a)).
Examples: bismuth, antimony, copper, silver, gold, quartz, water, air, hydro-
gen.

Note that diamagnetism occurs in all materials. However, the weak

diamagnetic force is suppressed in paramagnetic and ferromagnetic materials.

Diamagnetic susceptibility is largely independent of temperature. For a per-

fect diamagnet, χ = −1. Example: Superconductors exhibiting Meissner effect,
i.e., the complete expulsion of magnetic flux from inside a superconductor when
it is cooled to below its critical temperature.

Figure 10.1: Magnetic field lines in the vicinity of (a) Diamagnetic, (b) param-
agnetic, and (c) ferromagnetic materials.

10.2.2 Paramagnetic materials

Paramagnetism arises from the dipole moments of unpaired electrons in a ma-
terial. These atomic dipole moments are randomly oriented and add vectorially
to zero.
The dipoles align with the field when placed in an external applied magnetic
field, and the substances acquire a net magnetic moment along the direction of
the field.

3
 10.3 Magnetic moment of an atom

Paramagnetic materials lose their magnetization as soon as the applied mag-

netic field is removed.
They are feebly attracted by a magnet, and field lines tend to pass through
the substance (as shown in Figure 10.1(b)).
Examples: aluminium, platinum, chromium, oxygen, manganese.

10.2.3 Ferromagnetic materials

In ferromagnetic materials, atoms have permanent net non-zero magnetic dipole
moments.
They remain magnetized even when the external field is removed.
They are strongly attracted by a magnet, and field lines tend to crowd into
the substance (as shown in Figure 10.1(c)).
Examples: iron, cobalt, nickel, gadolinium, dysprosium, permalloy.

Ferromagnetic materials possess all the characteristics of paramagnetic ma-

terials with much greater intensity.

10.3 Magnetic moment of an atom

Figure 10.2: Magnetic dipole moment (µ = I × A) of an orbiting electron.

An electron circulating around the nucleus is equivalent to a magnetic dipole

(as schematically represented in Figure 10.2). Equivalent current of an electron

4
 10.3 Magnetic moment of an atom

of charge e revolving with time period T ,

e
i= (10.7)
T
Now,
2πr
T = (10.8)
v
where r is the radius of the orbit, and v the speed of the electron.
Therefore, from (10.7) and (10.8),
ev
i= (10.9)
2πr
The magnetic dipole moment associated with the electron

µm = i × A (current × area enclosed by the current loop)

ev
= × πr2
2πr
evr
= (10.10)
2
Now, we know from Bohr’s postulates that the angular momentum of the
electron is given by
nh
mvr = , n = 1, 2, 3...
2π
nh
or, vr = (10.11)
2πm
Therefore, from (10.10) and (10.11),
e nh
µm =×
2  2πm
eh
=⇒ µm = n (10.12)
4πm

 Thatis, the dipole moment is given by integral multiples of

eh
= 9.28 × 10−24 Am2 which is called the Bohr magneton.
4πm

5
 10.4 Paramagnetic susceptibility

10.4 Paramagnetic susceptibility

An expression for paramagnetic susceptibility can be obtained from Langevin’s
theory based on kinetic theory of gases. It is based on the notion that atomic
magnetic dipoles tend to align along an external magnetic field, but remain at
an angle with the field due to thermal agitation.
It can be shown that the paramagnetic susceptibility is given by

µ0 nµ2m
χ= (10.13)
3kT
where µm is the atomic magnetic dipole moment.

(10.13) can be rewritten as

C
χ= (10.14)
T
where
µ0 nµ2m
C= (10.15)
3k
is the Curie constant. It is material-specific.
Therefore, for paramagnetic materials,
1
χ∝ (10.16)
T
The variation of χ and 1/χ with T is shown in Figure 10.3. This is valid for
high T and low fields.

Figure 10.3: Temperature dependence of paramagnetic susceptibility.

6
 10.5 Ferromagnetism

10.5 Ferromagnetism
10.5.1 Domain theory
Pierre-Ernest Weiss proposed a hypothetical concept of ferromagnetic domains.
Some of the assumptions in the domains concept are as follows (illustrated in
Figure 10.4):

1. A ferromagnetic specimen contains a number of domains which are spon-

taneously magnetized.

2. Domains arise from certain mutual exchange interaction of neighbouring

atoms.

3. In bulk materials, domains are randomly oriented and cancel out each
other.

4. In presence of an external field, domain walls move to expand the favourable

domains.

5. Also, domains rotate and reorient in the direction of the applied field.

6. After the removal of the external field, domains may not completely revert
to the original random distribution.

Figure 10.4: Ferromagnetic domains.

7
 10.5 Ferromagnetism

10.5.2 Ferromagnetic susceptibility

To obtain an expression for ferromagnetic susceptibility, the applied field is
modified to the effective field, given by

Heff = H + λM (10.17)

where λ is a constant, called Weiss constant. λM is the internal molecular field.

It can be shown that ferromagnetic susceptibility is given by

µ0 nµ2m
χ= (10.18)
3k(T − Tc )
where
µ0 nµ2m λ
Tc = = Cλ (10.19)
3k
is called the Curie temperature. Curie temperature is the temperature above
which ferromagnetic materials lose their permanent magnetization and exhibit
paramagnetic behaviour. So, ferromagnetic susceptibility

Figure 10.5: Temperature dependence of ferromagnetic susceptibility.

C
χ= (10.20)
(T − Tc )
This is known as Curie-Weiss law of ferromagnetism.

8
 10.5 Ferromagnetism

More accurately (as schematically represented in Figure 10.5),

C
χ= (10.21)
(T − Tc )γ
for T > Tc , with γ ∼ 1. And
C
χ= (10.22)
(T − Tc )β
1
T < Tc , with β ∼ for T << Tc .
2

10.5.3 Domains and hysteresis

Once magnetic domains are reoriented by an external field, it takes some en-
ergy to change their orientation again. B does not vary linearly with H, and
consequently µr does not have a fixed value. In this case, the B(H) dependence
can be represented by a B − H curve or hysteresis curve.

Figure 10.6: A typical hysteresis curve.

A hysteresis curve or hysteresis loop is a B − H curve under the influence

of an AC magnetizing force, obtained by cycling the H field between extreme

9
 10.5 Ferromagnetism

positive and negative values. A typical hysteresis loop is shown in Figure 10.6
and described below.
As H is varied through (0 → +Hm → −Hm → +Hm ), the B − H curve first
follows the path from 0 to a, and then traces the closed loop abcdef a.
Subsequent field cycling through +Hm → −Hm → +Hm retraces the same
closed loop.
Increasing the applied field beyond ±Hm leaves the value of B fixed at the
saturation value of ±Bm (points a and d on the curve).
The non-zero value of B (given by ±Br (points b and e, respectively)) when
the applied field is reduced to zero from ±Hm is called residual magnetism or
remanence or retentivity.
The value of H at which B eventually goes to zero (points c and f for the
(+Hm → −Hm ) and (−Hm → +Hm ) cycles, respectively) is called the coercive
field, coercive force, or Coercivity, Hc .

10.5.4 Hard and soft ferromagnetic materials

Figure 10.7: Hard and soft ferromagnets.

Materials with large coercivity are called hard ferromagnets. Those with small
coercivity are called soft ferromagnets.
Hard ferromagnets with high retentivity are useful for construction of per-
manent magnets. Examples: Steel, AlNiCo, NdFeB, SmCo.
Soft ferromagnets are useful for making cores of electromagnets. Examples:
Soft iron, Permalloy.

Other characteristics crucial for making cores of electromagnets are:

10
 10.5 Ferromagnetism

1. high retentivity

2. low hysteresis loss (i.e., small area of the hysteresis loop)

3. high initial permeability.

********************* The End (of this Unit) *********************

11
 Mechanical properties of matter
PH-1007 (Physics)
Dr. Gorky Shaw

11.1 Elasticity and Hooke’s law

The resistance force acting per unit cross section area of a body subjected to
an external load is known as stress.
Mathematically,
Applied load
Stress = (11.1)
Cross section area
Stress has the unit of N/m2 (SI) or dyne/cm2 (CGS).
The deformation in the body due to the application of the external load is
given by strain.
Mathematically, strain is equal to the ratio of the change in a dimension (like
length) to its original value.
Strain is a ratio and is unit-less.
Stress versus strain variation in a typical ductile material is represented in
Figure 11.1.

Figure 11.1: Stress-strain curve for ductile materials.

 11.1 Elasticity and Hooke’s law
Some important features of the stress-strain curve
1. Proportional limit (OA): In this region, the stress-strain curve is a
straight line. That is, stress is directly proportional to strain and Hooke’s
law is obeyed.
2. Elastic limit (OB): It is the limit beyond which the material will no
longer go back to its original shape when the external force is removed.
In the regime AB the material still has elastic properties, which disappear
beyond the point B.
3. Upper yield point (B): is the point at which plastic deformation of the
material is initiated. At this point, maximum external load or stress is
required to initiate plastic deformation. The stress corresponding to B is
known as yield strength.
4. Lower yield point (C): At this point, minimum load is required to
maintain the plastic behaviour of the material.
5. Ultimate tensile strength (D): corresponds to the maximum stress
that the material can withstand before breaking. Beyond this point, neck-
ing starts inside the material.
6. Fracture / Breaking strength (E): is the point at which the material
fractures / breaks. The corresponding stress is known as breaking strength.
For comparison, typical stress-strain curves for brittle, ductile, and plastic ma-
terials are shown in Figure 11.2.

Figure 11.2: Stress-strain curves for brittle, ductile, and plastic materials.

2
 11.2 Elastic moduli

Hooke’s law : Within the elastic limit, stress is proportional to the corre-
sponding strain. i.e.,

stress ∝ strain
or, stress = E × strain
stress
or, E = (11.2)
strain
where the material-specific constant of proportionality E is known as elastic
modulus or modulus of elasticity and has the same unit as stress, i.e., N/m2
(SI) or dyne/cm2 (CGS).

11.2 Elastic moduli

There are three types elastic moduli corresponding to three types of strain:

1. Young’s modulus (E or Y ) corresponding to longitudinal strain.

2. Shear modulus of modulus rigidity (η) corresponding to shear strain.

3. Bulk modulus (K) corresponding to volume strain.

All elastic moduli have positive values by definition.

11.2.1 Young’s modulus (Y )

Figure 11.3: Young’s modulus.

longitudinal (or normal) stress

Young’s modulus =
longitudinal (or normal) strain
Fn /A Fn · l
In Figure 11.3, Y = = (11.3)
δl/l δl · A

3
 11.2 Elastic moduli

11.2.2 Shear modulus (η)

Figure 11.4: Modulus of rigidity.

shear stress
Shear modulus =
Shear strain
Ft /A Ft
In Figure 11.4, η = = (11.4)
δt/l A·γ

11.2.3 Bulk modulus (K)

Figure 11.5: Bulk modulus.

volume (or bulk) stress

Bulk modulus =
volume (or bulk) strain
P P ·V
In Figure 11.5, K = = (11.5)
−δV /V δV
where P is the pressure applied on the body. The negative sign indicates
that increasing pressure leads to decrease of volume.

4
 11.2 Elastic moduli

11.2.4 Poisson’s ratio

Figure 11.6: Poisson’s ratio.

The three types of stresses and strains are not independent of each other. For
example, a longitudinal stress producing longitudinal elongation also produces
lateral compression, so that the volume of the material is preserved (as shown
in Figure 11.6).
Poisson’s ratio (σ) is a measure of the Poisson effect that describes the ex-
pansion or contraction of a material in directions perpendicular to the direction
of loading. It is the ratio between longitudinal strain (denoted by α) and lateral
strain (denote by β).
Mathematically,
lateral strain β
Poisson’s ratio = − =− (11.6)
longitudinal strain α
In Figure 11.6,
δw δd
β=− =−
w d
δl
and, α = (11.7)
l
Therefore, Poisson’s ratio
β δw/w δd/d
σ=− = = (11.8)
α δl/l δl/l
1
Poisson’s ratio has no unit and has values lying between + and −1.
2

5
 11.3 Relation between elastic constants

11.3 Relation between elastic constants

The three elastic moduli Y, K, η and Poisson’s ratio σ are related to each other.

11.3.1 Relation between Y, K and σ

Derivation not required for examination.

We consider a force F acting normally outwards on each of the six faces of a

cube of unit length. Thus, a uniform deformation is produced in all directions
in the cube.
Now, we consider the force acting along the X-axis. If this force produces a
longitudinal extension α along the X-axis, then we have 1
F
Y = (11.9)
α
The same force also produces a lateral strain in Y and Z directions. If lateral
contraction is β, then we have 2
β
σ= (11.10)
α
Hence, from (11.9) and (11.10),
F
β = σα = σ (11.11)
Y
Similarly, the force along Y-axis produces an extension in Y-direction and
contractions in X and Z directions. And the force along Z-axis produces an
extension along Z-direction and contractions along X and Y directions.
Therefore, net change in length along each axis
F F F
e = α − 2β = − 2σ = (1 − 2σ) (11.12)
Y Y Y
Hence, the change in volume of the cube is

δV = (1 + e)3 − 1 ' 3e (11.13)

ignoring higher order terms in e, assuming e << 1.

1
The numerator of (11.9) is just F because area of the cube is unity.
2
The minus sign due to lateral contraction is absorbed in β itself.

6
 11.3 Relation between elastic constants

Hence, bulk modulus of the material is given by

F F FY
K= = = (11.14)
δV 3e 3F (1 − 2σ)
or,

Y = 3K(1 − 2σ) (11.15)

11.3.2 Relation between Y, η and σ

Derivation not required for examination.

We consider a compression stress F acting along X-direction and an extension

stress F acting along Y-direction on a cube of unit length.
The force along X-axis produces a longitudinal compression α along X-
direction and a lateral extension β along Y-direction. Similarly, the force along
Y-axis produces a longitudinal extension α along Y-direction and a lateral com-
pression β along X-direction.
Hence, net change in length along the X and Y directions, respectively, are
given by
F F F
ex = −α − β = − − σ = − (1 + σ)
Y Y Y
F
and, ey = α + β = (1 + σ) (11.16)
Y
Also,

ez = 0 (11.17)

because along Z-direction, the compressive stress along X-axis produces lateral
extension, and the extensive stress along Y-axis produces equal lateral com-
pression.
It can be shown that a force F of compression along X-axis and an equal
force of extension along Y-axis are equivalent to a shear stress of magnitude F
at 450 to X and Y axes. This shear stress produces a shearing strain of angle θ
in the XY-plane, such that

θ = 2ex = 2ey (in magnitude) (11.18)

7
 11.3 Relation between elastic constants

Hence, shear modulus is given by

F F FY
η= = = (11.19)
θ 2ex 2F (1 + σ)
Therefore,

Y = 2η(1 + σ) (11.20)

11.3.3 Relation between K, η and σ

From (11.15) and (11.20), we can write

2η(1 + σ) = 3K(1 − 2σ)

=⇒ (6K + 2η)σ = 3K − 2η
(11.21)

Therefore,

3K − 2η
σ= (11.22)
6K + 2η

11.3.4 Relation between Y, K and η

Using (11.22) in (11.20), we get
 
3K − 2η
Y = 2η 1 +
6K + 2η
η
= (6K + 2η + 3K − 2η) (11.23)
3K + η
Therefore,

9Kη
Y = (11.24)
3K + η

8
 11.3 Relation between elastic constants

11.3.5 Limiting values of Poisson’s ratio (σ)

From (11.15), (11.20) and (11.21),

Y = 2η(1 + σ) = 3K(1 − 2σ) (11.25)

Now, by definition Y > 0, η > 0 and K > 0. Therefore, we have

(1 + σ) > 0 and (1 − 2σ) > 0

i.e., σ > −1 and 1 > 2σ
1
i.e., − 1 < σ and σ <
2
(11.26)

Therefore, we have,
1
−1 < σ < (11.27)
2

******************************** FIN ********************************