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Lectures On Quantum Theory Complete

This document provides lecture notes for a course on quantum theory taught by Dr. Frederic P. Schuller in 2015. The notes were taken and later expanded upon by Simon Rea and Richie Dadhley. The course covered fundamental topics in quantum theory including axioms, Hilbert spaces, operators, spectra, spin, harmonic oscillators, and measurements. The notes are made freely available online through the authors' blogs.

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0% found this document useful (0 votes)

443 views

Lectures On Quantum Theory Complete

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Allan Pineda

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Download as PDF, TXT or read online on Scribd

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Lectures on Quantum Theory

Course delivered in 2015 by

Dr Frederic P. Schuller
Friedrich-Alexander-Universität Erlangen-Nürnberg,
Institut für Theoretische Physik III

E-mail: [email protected]

Notes taken by Simon Rea & Richie Dadhley

E-mail: [email protected]
E-mail: [email protected]

Last updated on January 31, 2019

Acknowledgments
This set of lecture notes accompanies Frederic Schuller’s course on Quantum Theory, taught
in the summer of 2015 at the Friedrich-Alexander-Universität Erlangen-Nürnberg as part
of the Elite Graduate Programme.
The entire course is hosted on YouTube at the following address:

www.youtube.com/playlist?list=PLPH7f_7ZlzxQVx5jRjbfRGEzWY_upS5K6

These lecture notes are not endorsed by Dr. Schuller or the University.
While I have tried to correct typos and errors made during the lectures (some helpfully
pointed out by YouTube commenters), I have also taken the liberty to add and/or modify
some of the material at various points in the notes. Any errors that may result from this
are, of course, mine.
If you have any comments regarding these notes, feel free to get in touch. Visit my
blog for the most up to date version of these notes

http://mathswithphysics.blogspot.com

My gratitude goes to Dr. Schuller for lecturing this course and for making it available
on YouTube.

Simon Rea

I have picked up the notes from where Simon left off (Lecture 8). I have also been
through the first lectures and added small details that I thought helpful. I have tried
to stay consistent with Simon’s inclusion of additional valuable material throughout the
remainder of the course. As with Simon’s comment above, any mistakes made because of
this are, of course, mine.
I have also made the up-to-date version of the notes available on my blog site

https://richie291.wixsite.com/theoreticalphysics

I would like to extend a message of thanks to Simon for providing these notes (and the
notes for Dr. Schuller’s other courses) online, I have personally found them very useful.
I would also like to show my gratitude to Dr. Schuller for putting his courses on
YouTube, I have found them both very informative and interesting, a credit to his brilliant
teaching ability.

Richie Dadhley
Contents

Introduction 1

1 Axioms of quantum mechanics 2

1.1 Desiderata 2
1.2 The axioms of quantum mechanics 5

2 Banach spaces 9
2.1 Generalities on Banach spaces 9
2.2 Bounded linear operators 13
2.3 Extension of bounded linear operators 21

3 Separable Hilbert Spaces 25

3.1 Relationship between norms and inner products 25
3.2 Hamel versus Schauder 33
3.3 Separable Hilbert spaces 35
3.4 Unitary maps 36

4 Projectors, bras and kets 41

4.1 Projectors 41
4.2 Closed linear subspaces 42
4.3 Orthogonal projections 43
4.4 Riesz representation theorem, bras and kets 46

5 Measure theory 50
5.1 General measure spaces and basic results 50
5.2 Borel σ-algebras 56
5.3 Lebesgue measure on Rd 58
5.4 Measurable maps 59
5.5 Push-forward of a measure 60

6 Integration of measurable functions 61

6.1 Characteristic and simple functions 61
6.2 Integration of non-negative measurable simple functions 63
6.3 Integration of non-negative measurable functions 65
6.4 Lebesgue integrable functions 71
6.5 The function spaces Lp (M, Σ, µ) 75

–i–
7 Self-adjoint and essentially self-adjoint operators 80
7.1 Adjoint operators 80
7.2 The adjoint of a symmetric operator 81
7.3 Closability, closure, closedness 82
7.4 Essentially self-adjoint operators 84
7.5 Criteria for self-adjointness and essential self-adjointness 85

8 Spectra and perturbation theory 87

8.1 Resolvent map and spectrum 87
8.2 The spectrum of a self-adjoint operator 88
8.3 Perturbation theory for point spectra of self-adjoint operators 89

9 Case study: momentum operator 95

9.1 The Momentum Operator 95
9.2 Absolutely Continuous Fucntions and Sobolev Spaces 95
9.3 Momentum Operator on a Compact Interval 97
9.4 Momentum Operator on a Circle 101

10 Inverse spectral theorem 103

10.1 Projection-valued measures 104
10.2 Real and Complex Valued Borel Measures Induced by a PVM 106
10.3 Integration With Respect to a PVM 106
10.4 The Inverse Spectral Theorem 112

11 Spectral theorem 113

11.1 Measurable Function Applied To A Spectrally Decomposable Self Adjoint
Operator 113
11.2 Reconstruct PVM From a Spectrally Decomposable, Self Adjoint Operator 114
11.3 Construction Of PVM From A Self Adjoint Operator 117
11.4 Commuting Operators 118

12 Stone’s Theorem and Construction of Observables 121

12.1 One Parameter Groups and Their Generators 121
12.2 Stone’s Theorem 123
12.3 Domains of Essential Self Adjointness ("Cores") 126
12.4 Position, Momentum and Angular Momentum 126
12.5 Schwartz Space S(Rd ) 128

13 Spin 130
13.1 General Spin 132
13.2 Derivation of Pure Point Spectrum 132
13.3 Pure Spin-j Systems 138

– ii –
14 Composite Systems 140
14.1 Tensor Product of Hilbert Spaces 141
14.2 Practical Rules for Tensor Products of Vectors 146
14.3 Tensor Product Between Operators 146
14.4 Symmetric and Antisymmetric Tensor Products 147
14.5 Collapse of Notation 149
14.6 Entanglement 149

15 Total Spin of Composite System 151

15.1 Total Spin 151
15.2 Eigenbasis For The Composite System in Terms of Simultaneous Eigenvec-
tors of A2 ⊗ 1, 1 ⊗ B 2 , A3 ⊗ 1 and 1 ⊗ B 3 152
15.3 Conversion to Eigenbasis in Terms of J , J3 , A ⊗ 1, 1 ⊗ B
2 2 2 153
15.4 Value of m 154
15.5 Clebsch-Gordan Coefficients for Maximal j 155
15.6 Clebsch-Gordan Coefficients For Lower Than Max j 155
15.7 Total Spin of Composite System 156

16 Quantum Harmonic Oscillator 158

16.1 What People Say 158
16.2 The Quantum Harmonic Oscillator 159
16.3 The Energy Spectrum 160

17 Measurements of Observables 165

17.1 Spectral Decomposition of H For The Quantum Harmonic Oscillator 165
17.2 Probability to Measure a Certain Value 167
17.3 Measurements of Pure States 168
17.4 Preparation of Pure States 171
17.5 Mixed States 172

18 The Fourier Operator 174

18.1 The Fourier Operator on Schwartz Space 174
18.2 Inverse of Fourier Operator 177
18.3 Extension of F to L2 (Rd ) 179
18.4 Convolutions 180

19 The Schrodinger Operator 182

19.1 Domain of Self Adjointness of Hfree 182
19.2 Spectrum of Hfree 183
19.3 Time Evolution 184
19.4 Physical Interpretation (for the Patient Observer) 185

– iii –
20 Periodic Potentials 187
20.1 Basics of Rigged Hilbert Space 188
20.2 Fundamental Solutions and Fundamental Matrix 188
20.3 Translation Operator 190
20.4 Application to Our Quantum Problem 192
20.5 Energy Bands 193
20.6 Quantitative Calculation (Outline) 193

21 Relativistic Quantum Mechanics 194

21.1 Heuristic Derivation of the Schrödinger Equation 194
21.2 Heuristic Derivation of the Relativistic Schrödinger Equation (Klein-Gordan
Equation) 195
21.3 Dirac Equation 196
21.4 The Deep Route of the Problem 197
21.5 Feynman Diagrams 198

Hence, we need to turn to new mathematics in which we can define a notion of

observable that allows for a spectrum of measurement results for a quantum observable
A of the form
σ(A) = discrete part ∪ continuous part.
An example would be the energies of the hydrogen atom

σ(H) =
−13.6 eV 0 eV

Note that one of the parts may actually be empty. For instance, as we will later show,
the simple quantum harmonic oscillator has the following energy spectrum

σ(H) =
1
2 ~ω ( 12 + n)~ω

while the spectrum of the position operator Q is σ(Q) = R.

2
Educated term for “wishlist”.

–2–
Also, the continuous part need not be connected, as is the case with spectrum of the
Hamiltonian an electron in a periodic potential

σ(H) =

It turns out that self-adjoint linear maps on a complex Hilbert space provide a suitable
formalism to describe the observables of quantum mechanics.

(b) An irreducible impact that each measurement has on the state of a quantum system.
The crucial example demonstrating this is the Stern-Gerlach experiment, which con-
sists in the following. Silver atoms are heated up in an oven and sent against a screen
with a hole. The atoms passing through the hole are then subjected to an inho-
mogeneous magnetic field, which deflects them according to the component of their
angular momentum in the direction of the field. Finally, a screen detects the various
deflections.

S Wτϕ

Ag
atoms

inhomogeneous
magnetic field

Since the angular momentum distribution of the silver atoms coming from the oven is
random, we would expect an even distribution of values of the component along the
direction of the magnetic field to be recorded on the final screen, as in S. However,
the impact pattern actually detected is that on the W τ ϕ screen. In fact, 50% of
the incoming atoms impact at the top and we say that their angular momentum
component is ↑, and the other 50% hit the bottom region, and we say that their
angular momentum component is ↓. This is another instance of our earlier point:
there seem to be only two possible values for the component of angular momentum
along the direction of the magnetic field, i.e. the spectrum is discrete. Hence, this is
not particularly surprising at this point.
Let us now consider successive iterations of this experiment. Introduce some system
of cartesian coordinates (x, y, z) and let SG(x) and SG(z) denotes a Stern-Gerlach
apparatus whose magnetic field points in the x and z-direction, respectively.

–3–
Suppose that we sent the atoms through a first SG(z) apparatus, and then we use
the z ↑ -output as the input of a second SG(z) apparatus.

z↑
100%
z↑
SG(z)
z↓
SG(z) 0%
z↓

The second SG(z) apparatus finds no z ↓ -atoms. This is not surprising since, intu-
itively, we “filtered out” all the z ↓ -atoms with the first apparatus. Suppose now that
we feed the z ↑ output of a SG(z) apparatus into a SG(x) apparatus.

x↑
50%
z↑
SG(x)
x↓
SG(z) 50%
z↓

Experimentally, we find that about half of the atoms are detected in the state x↑ and
half in the state x↓ . This is, again, not surprising since we only filtered out the z ↑
atoms, and hence we can interpret this result as saying that the x↑ , x↓ states are
independent from the z ↑ , z ↓ .
If our ideas of “filtering states out” is correct, then feeding the x↑ -output of the
previous set-up to another SG(z) apparatus should clearly produce a 100% z ↑ -output,
since we already filtered out all the z ↓ ones in the previous step.

z↑
50%
x↑
SG(z)
z↑ z↓
SG(x) 50%
x↓
SG(z)
z↓

Surprisingly, the output is again 50-50. The idea behind this result is the following.
The SG(z) apparatus left the atoms in a state such that a repeated measurement
with the SG(z) apparatus would give the same result, and similarly for the SG(x)

–4–
apparatus. However, the measurement of the SG(x) apparatus somehow altered the
state of the atoms in such a way as to “reset” them with respect to a measurement
by the SG(z) apparatus. For more details on the Stern-Gerlach experiment and
further conclusions one can draw from its results, you should consult the book Modern
Quantum Mechanics by J. J. Sakurai. The conclusion that we are interested in here
is that measurements can alter the state of a system.

(c) Even if the state ρ of a quantum system is completely known, the only prediction one
can make for the measurement of some observable A is the probability that the mea-
sured value, which is an element of the spectrum σ(A), lies within a Borel-measurable
subset E ⊆ R, denoted by µA ρ (E).

In particular, one cannot predict which concrete outcome an isolated measurement

will produce. This is even more annoying given that the precise impact that a mea-
surement has on the state of the system (see previous point) depends on the observed
outcome of the measurement.

A suitable theory that accommodates all known experimental facts has been developed
between 1900 and 1927 on the physics side by, among others, Schrödinger, Heisenberg and
Dirac, and on the mathematical side almost single-handedly by von Neumann who invented
a massive proportion of a field known today as functional analysis.

1.2 The axioms of quantum mechanics

We will now present the axioms of quantum mechanics by using notions and terminology
that will be defined later in the course. In this sense, this section constitutes a preview of
the next few lectures.

Axiom 1 (Quantum systems and states). To every quantum system there is asso-
ciated a separable complex Hilbert space (H, +, ·, h·|·i). The states of the system are
all positive, trace-class linear maps ρ : H → H for which Tr ρ = 1.

Remark 1.1 . Throughout the quantum mechanics literature, it is stated that the unit, or
normalised, elements ψ ∈ H (that is, hψ|ψi = 1) are the states of the quantum system.
This is not correct.
States can be pure or mixed. A state ρ : H → H is called pure if

hψ|αi
∃ ψ ∈ H : ∀ α ∈ H : ρ(α) = ψ.
hψ|ψi

Thus, we can associate to each pure state ρ an element ψ ∈ H. However, this correspondence
is not one-to-one. Even if we restrict to pure states and impose the normalisation condition,
there can be many ψ ∈ H representing the same pure state ρ.
Therefore, it is wrong to say that the states of the quantum system are the normalised
elements of the Hilbert space, since they do not represent all the states of the system, and
do not even represent uniquely the states that they do represent.

–5–
The terms used in Axiom 1 are defined as follows.

Definition. A complex Hilbert space is a tuple (H, +, ·, h·|·i) where

• H is a set

• + is a map + : H × H → H

• · is a map · : C × H → H (typically suppressed in the notation)

such that the triple (H, +, ·) is a vector space over C, and

• h·|·i is a sesqui-linear3 inner product, i.e. a map h·|·i : H × H → C satisfying

(i) hϕ|ψi = hψ|ϕi (conjugate symmetry/Hermitian property)

for all ϕ, ψ1 , ψ2 ∈ H and z ∈ C,

and moreover

• H is a complete metric space with respect to the metric induced by the norm induced
in turn by the sesqui-linear map h·|·i. Explicitly, for every sequence φ : N → H that
satisfies the Cauchy property, namely

∀ ε > 0 : ∃ N ∈ N : ∀ n, m ≥ N : kφn − φm k < ε,

where φn := φ(n) and kψk := hψ|ψi, then the sequence converges in H, i.e.
p

∃ ϕ ∈ H : ∀ ε > 0 : ∃ N ∈ N : ∀ n ≥ N : kϕ − φn k < ε.

Note that the C-vector space (H, +, ·) need not be finite-dimensional and, in fact, we
will mostly work with infinite-dimensional Hilbert spaces.

Definition. A map A : DA → H, where the subspace DA ⊆ H is called the domain of A,

is a linear map if

∀ ϕ, ψ ∈ DA : ∀ z ∈ C : A(zϕ + ψ) = zA(ϕ) + A(ψ).

From now on, if there is no risk of confusion, we will write Aϕ := A(ϕ) in order to
spare some brackets. We will be particularly interested in special types of linear map.

Definition. A linear map A : DA → H is densely defined if DA is dense in H, i.e.

∀ ψ ∈ H : ∀ ε > 0 : ∃ α ∈ DA : kα − ψk < ε.

Definition. A linear map A : DA → H is said to be positive if

∀ ψ ∈ DA : hψ|Aψi ≥ 0.
3
sesqui is Latin for “one and a half”.

–6–
Definition. A linear map A : DA → H is said to be of trace-class if DA = H and, for any
orthonormal basis {en } of H, the sum/series
X
hen |Aen i < ∞.
n

If A : H → H is of trace-class, one can show that the value of does not

P
n hen |Aen i
depend on the choice of orthonormal basis {en }.
Definition. Let A : H → H be of trace-class. Then the trace of A is
X
Tr A := hen |Aen i
n

where {en } is any orthonormal basis of H.

Axiom 2 (Observables). The observables of a quantum system are the self-adjoint

linear maps A : DA → H.

While the notion of a self-adjoint map is easy to define in finite-dimensional spaces, it

is much more subtle for infinite-dimensional spaces.
Definition. A densely defined linear map A : DA → H is said to be of self-adjoint if it
coincides with its adjoint map A∗ : DA∗ → H, that is
• DA = DA∗

• ∀ ϕ ∈ DA : Aϕ = A∗ ϕ.
Definition. The adjoint map A∗ : DA∗ → H of a linear map A : DA → H is defined by
• DA∗ := {ψ ∈ H | ∀ α ∈ DA : ∃ η ∈ H : hψ|Aαi = hη|αi}

• A∗ ψ := η.
We will later show that the adjoint map is well-defined, i.e. for each α ∈ DA and ψ ∈ H
there exists at most one η ∈ H such that hψ|Aαi = hη|αi.
Remark 1.2 . If we defined DA∗ by requiring that η ∈ DA , we would obtain a notion of self-
adjointness which has undesirable properties. In particular, the spectrum (to be defined
later) of a self-adjoint operator would not be guaranteed to be a subset of R.

Axiom 3 (Measurement). The probability that a measurement of an observable A

on a system that is in the state ρ yields a result in the Borel set E ⊆ R is given by

µA
ρ (E) := Tr(PA (E) ◦ ρ)

where the map PA : Borel(R) → L(H), from the Borel-measurable subsets of R to the
Banach space of bounded linear maps on H, is the unique projection-valued measure
that is associated with the self-adjoint map A according to the spectral theorem.

–7–
We will later see that the composition of a bounded linear map with a trace-class map
is again of trace-class, so that Tr(PA (E) ◦ ρ) is well-defined. For completeness, the spectral
theorem states that for any self-adjoint map A there exists a projection-valued measure PA
such that A can be represented in terms of the Lebesgue-Stieltjes integral as
Z
A= λ dPA (λ).
R

This is the infinite-dimensional analogue of the diagonalisation theorem for symmetric or

Hermitian matrices on finite-dimensional vector spaces, and it is the theorem in which the
first half of the course will find its climax.

Axiom 4 (Unitary dynamics). In a time interval (t1 , t2 ) ⊆ R in which no measure-

ment occurs, the state ρ at time t1 , denoted ρ(t1 ), is related to the state ρ at time t2 ,
denoted ρ(t2 ), by
ρ(t2 ) = U (t2 − t1 )ρ(t1 )U −1 (t2 − t1 )
with the unitary evolution operator U defined as

U (t) := exp − ~i Ht ,

where H is the energy observable and, for any observable A and f : R → C, we define
Z
f (A) := f (λ) dPA (λ).
R

Note that, as was the case for the previous axiom, the spectral theorem is crucial since
it is needed to define the unitary evolution operator.

Axiom 5 (Projective dynamics). The state ρafter of a quantum system immediately

following the measurement of an observable A is

PA (E) ◦ ρbefore ◦ PA (E)

ρafter :=
Tr(PA (E) ◦ ρbefore ◦ PA (E))

where ρbefore is the state immediately preceding the measurement and E ⊆ R is the
smallest Borel set in which the actual outcome of the measurement happened to lie.

–8–
2 Banach spaces
Hilbert spaces are a special type of a more general class of spaces known as Banach spaces.
We are interested in Banach spaces not just for the sake generality, but also because they
naturally appear in Hilbert space theory. For instance, the space of bounded linear maps
on a Hilbert space is not itself a Hilbert space, but only a Banach space.

2.1 Generalities on Banach spaces

We begin with some basis notions from metric space theory.
Definition. A metric space is a pair (X, d), where X is a set and d is a metric on X, that
is, a map d : X × X → R satisfying
(i) d(x, x) ≥ 0 (non-negativity)

(ii) d(x, y) = 0 ⇔ x = y (identity of indiscernibles)

(iii) d(x, y) = d(y, x) (symmetry)

(iv) d(x, z) ≤ d(x, y) + d(y, z) (triangle inequality)

for all x, y, z ∈ X.
Definition. A sequence {xn }n∈N in a metric space (X, d) is said to converge to an element
x ∈ X, written lim xn = x, if
n→∞

∀ ε > 0 : ∃ N ∈ N : ∀ n ≥ N : d(xn , x) < ε.

A sequence in a metric space can converge to at most one element.
Definition. A Cauchy sequence in a metric space (X, d) is a sequence {xn }n∈N such that
∀ ε > 0 : ∃ N ∈ N : ∀ n, m ≥ N : d(xn , xm ) < ε.
Any convergent sequence is clearly a Cauchy sequence.
Definition. A metric space (X, d) is said to be complete if every Cauchy sequence converges
to some x ∈ X.

A natural metric on a vector space is that induced by a norm.

Definition. A normed space is a (complex) vector space (V, +, ·) equipped with a norm,
that is, a map k · k : V → R satisfying

(i) kf k ≥ 0 (non-negativity)

(ii) kf k = 0 ⇔ f = 0 (definiteness)

(iii) kz · f k = |z|kf k (homogeneity/scalability)

(iv) kf + gk ≤ kf k + kgk (triangle inequality/sub-additivity)

for all f, g ∈ V and all z ∈ C.

–9–
Once we have a norm k · k on V , we can define a metric d on V by

d(f, g) := kf − gk.

Then we say that the normed space (V, k · k) is complete if the metric space (V, d), where
d is the metric induced by k · k, is complete. Note that we will usually suppress inessential
information in the notation, for example writing (V, k · k) instead of (V, +, ·, k · k).

Definition. A Banach space is a complete normed vector space.

Example 2.1 . The space CC0 [0, 1] := {f : [0, 1] → C | f is continuous}, where the continuity
is with respect to the standard topologies on [0, 1] ⊂ R and C, is a Banach space. Let us
show this in some detail.

Proof. (a) First, define two operations +, · pointwise, that is, for any x ∈ [0, 1]

(f + g)(x) := f (x) + g(x) (z · f )(x) := zf (x).

Suppose that f, g ∈ CC0 [0, 1], that is

∀ x0 ∈ [0, 1] : ∀ ε > 0 : ∃ δ > 0 : ∀ x ∈ (x0 − δ, x0 + δ) : |f (x) − f (x0 )| < ε

and similarly for g. Fix x0 ∈ [0, 1] and ε > 0. Then, there exist δ1 , δ2 > 0 such that
ε
∀ x ∈ (x0 − δ1 , x0 + δ1 ) : |f (x) − f (x0 )| < 2
ε
∀ x ∈ (x0 − δ2 , x0 + δ2 ) : |g(x) − g(x0 )| < 2.

Let δ := min{δ1 , δ2 }. Then, for all x ∈ (x0 − δ, x0 + δ), we have

|(f + g)(x) − (f + g)(x0 )| := |f (x) + g(x) − (f (x0 ) + g(x0 ))|

= |f (x) − f (x0 ) + g(x) − g(x0 ))|
≤ |f (x) − f (x0 )| + |g(x) − g(x0 ))|
ε ε
< 2 + 2
= ε.

Since x0 ∈ [0, 1] was arbitrary, we have f + g ∈ CC0 [0, 1]. Similarly, for any z ∈ C and
f ∈ CC0 [0, 1], we also have z · f ∈ CC0 [0, 1]. It is immediate to check that the complex
vector space structure of C implies that the operations

+ : CC0 [0, 1] × CC0 [0, 1] → CC0 [0, 1] · : C × CC0 [0, 1] → CC0 [0, 1]
(f, g) 7→ f + g (z, f ) 7→ z · f

make (CC0 [0, 1], +, ·) into a complex vector space.

(b) Since [0, 1] is closed and bounded, it is compact and hence every complex-valued
continuous function f : [0, 1] → C is bounded, in the sense that

sup |f (x)| < ∞.

x∈[0,1]

– 10 –
We can thus define a norm on CC0 [0, 1], called the supremum (or infinity) norm, by

kf k∞ := sup |f (x)|.
x∈[0,1]

Let us show that this is indeed a norm on (CC0 [0, 1], +, ·) by checking that the four
defining properties hold. Let f, g ∈ CC0 [0, 1] and z ∈ C. Then

(b.i) kf k∞ := sup |f (x)| ≥ 0 since |f (x)| ≥ 0 for all x ∈ [0, 1].

x∈[0,1]
(b.ii) kf k∞ = 0 ⇔ sup |f (x)| = 0. By definition of supremum, we have
x∈[0,1]

∀ x ∈ [0, 1] : |f (x)| ≤ sup |f (x)| = 0.

x∈[0,1]

But since we also have |f (x)| ≥ 0 for all x ∈ [0, 1], f is identically zero.
(b.iii) kz · f k∞ := sup |zf (x)| = sup |z||f (x)| = |z| sup |f (x)| = |z|kf k∞ .
x∈[0,1] x∈[0,1] x∈[0,1]
(b.iv) By using the triangle inequality for the modulus of complex numbers, we have

kf + gk∞ := sup |(f + g)(x)|

x∈[0,1]
= sup |f (x) + g(x)|
x∈[0,1]
≤ sup (|f (x)| + |g(x)|)
x∈[0,1]
= sup |f (x)| + sup |g(x)|
x∈[0,1] x∈[0,1]
= kf k∞ + kgk∞ .

Hence, (CC0 [0, 1], k · k∞ ) is indeed a normed space.

∀ ε > 0 : ∃ N ∈ N : ∀ n, m ≥ N : kfn − fm k∞ < ε.

We seek an f ∈ CC0 [0, 1] such that lim fn = f . We will proceed in three steps.
n→∞

(c.i) Fix y ∈ [0, 1] and ε > 0. By definition of supremum, we have

|fn (y) − fm (y)| ≤ sup |fn (x) − fm (x)| =: kfn − fm k∞ .

x∈[0,1]

Hence, there exists N ∈ N such that

∀ n, m ≥ N : |fn (y) − fm (y)| < ε,

that is, the sequence of complex numbers {fn (y)}n∈N is a Cauchy sequence. Since
C is a complete metric space4 , there exists zy ∈ C such that lim fn (y) = zy .
n→∞
4
The standard metric on C is induced by the modulus of complex numbers.

– 11 –
Thus, we can define a function

f : [0, 1] → C
x 7→ zx ,

called the pointwise limit of f , which by definition satisfies

∀ x ∈ [0, 1] : lim fn (x) = f (x).

n→∞

Note that this does not automatically imply that lim fn = f (converge with
n→∞
respect to the supremum norm), nor that f ∈ CC0 [0, 1], and hence we need to
check separately that these do, in fact, hold.
(c.ii) First, let us check that f ∈ CC0 [0, 1], that is, f is continuous. Let x0 ∈ [0, 1] and
ε > 0. For each x ∈ [0, 1], we have

|f (x) − f (x0 )| = |f (x) − fn (x) + fn (x) − fn (x0 ) + fn (x0 ) − f (x0 )|

≤ |f (x) − fn (x)| + |fn (x) − fn (x0 )| + |fn (x0 ) − f (x0 )|.

Since f is the pointwise limit of {fn }n∈N , for each x ∈ [0, 1] there exists N ∈ N
such that
∀ n ≥ N : |f (x) − fn (x)| < 3ε .
In particular, we also have

∀ n ≥ N : |fn (x0 ) − f (x0 )| < 3ε .

Moreover, since fn ∈ CC0 [0, 1] by assumption, there exists δ > 0 such that

∀ x ∈ (x0 − δ, x0 + δ) : |fn (x) − fn (x0 )| < 3ε .

Fix n ≥ N . Then, it follows that for all x ∈ (x0 − δ, x0 + δ), we have

Since x0 ∈ [0, 1] was arbitrary, we have f ∈ CC0 [0, 1].

(c.iii) Finally, it remains to show that lim fn = f . To that end, let ε > 0. By the
n→∞
triangle inequality for k · k∞ , we have

kfn − f k∞ = kfn − fm + fm − f k∞
≤ kfn − fm k∞ + kfm − f k∞ .

Since {fn }n∈N is Cauchy by assumption, there exists N1 ∈ N such that

∀ n, m ≥ N1 : kfn − fm k∞ < 2ε .

– 12 –
Moreover, since f is the pointwise limit of {fn }n∈N , for each x ∈ [0, 1] there
exists N2 ∈ N such that
∀ m ≥ N2 : |fm (x) − f (x)| < 2ε .
By definition of supremum, we have
∀ m ≥ N2 : kfm − f k∞ = sup |fm (x) − f (x)| ≤ 2ε .
x∈[0,1]

Let N := max{N1 , N2 } and fix m ≥ N . Then, for all n ≥ N , we have

ε ε
kfn − f k∞ ≤ kfn − fm k∞ + kfm − f k∞ < 2 + 2 = ε.
Thus, lim fn = f and we call f the uniform limit of {fn }n∈N .
n→∞
This completes the proof that (CC0 [0, 1], k · k∞ ) is a Banach space.

Remark 2.2 . The previous example shows that checking that something is a Banach space,
and the completeness property in particular, can be quite tedious. However, in the following,
we will typically already be working with a Banach (or Hilbert) space and hence, rather
than having to check that the completeness property holds, we will instead be able to use
it to infer the existence (within that space) of the limit of any Cauchy sequence.

2.2 Bounded linear operators

As usual in mathematics, once we introduce a new types of structure, we also want study
maps between instances of those structures, with extra emphasis placed on the structure-
preserving maps. We begin with linear maps from a normed space to a Banach space.
Definition. Let (V, k · kV ) be a normed space and (W, k · kW ) a Banach space. A linear
map, also called a linear operator, A : V → W is said to be bounded if
kAf kW
sup < ∞.
f ∈V kf kV
Note that the quotient is not defined for f = 0. Hence, to be precise, we should write
V \ {0} instead of just V . Let us agree that is what mean in the above definition. There
are several equivalent characterisations of the boundedness property.
Proposition 2.3. A linear operator A : V → W is bounded if, and only if, any of the
following conditions are satisfied.

(i) sup kAf kW < ∞

kf kV =1

(ii) ∃ k > 0 : ∀ f ∈ V : kf kV ≤ 1 ⇒ kAf kW ≤ k

(iii) ∃ k > 0 : ∀ f ∈ V : kAf kW ≤ kkf kV

(iv) the map A : V → W is continuous with respect to the topologies induced by the respec-
tive norms on V and W

(v) the map A is continuous at 0 ∈ V .

– 13 –
The first one of these follows immediately from the homogeneity of the norm. Indeed,
suppose that kf kV 6= 1. Then

kAf kW
= kf k−1 −1
V kAf kW = kA(kf kV f )kW = kAf kW
e
kf kV

where fe := kf k−1
V f is such that kf kV = 1. Hence, the boundedness property is equivalent
e
to condition (i) above.

Definition. Let A : V → W be a bounded operator. The operator norm of A is defined as

kAf kW
kAk := sup kAf kW = sup .
kf kV =1 f ∈V kf kV

Example 2.4 . Let idW : W → W be the identity operator on a Banach space W . Then

k idW f kW
sup = sup 1 = 1 < ∞.
f ∈W kf kW f ∈W

Hence, idW is a bounded operator and has unit norm.

Example 2.5 . Denote by CC1 [0, 1] the complex vector space of once continuously differen-
tiable complex-valued functions on [0, 1]. Since differentiability implies continuity, this is
a subset of CC0 [0, 1]. Moreover, since sums and scaling by a complex number of continu-
ously differentiable functions are again continuously differentiable, this is, in fact, a vector
subspace of CC0 [0, 1], and hence also a normed space with the supremum norm k · k∞ .
Consider the first derivative operator

D : CC1 [0, 1] → CC0 [0, 1]

f 7→ f 0 .

We know from undergraduate real analysis that D is a linear operator. We will now show
that D is an unbounded5 linear operator. That is,
kDf k∞
sup = ∞.
1 [0,1]
f ∈CC kf k∞

Note that, since the norm is a function into the real numbers, both kDf k∞ and kf k∞ are
always finite for any f ∈ CC1 [0, 1]. Recall that the supremum of a set of real numbers is its
least upper bound and, in particular, it need not be an element of the set itself. What we
have to show is that the set

kDf k∞ 1
f ∈ CC [0, 1] ⊂ R
kf k∞
contains arbitrarily large elements. One way to do this is to exhibit a positively divergent
(or unbounded from above) sequence within the set.
5
Some people take the term unbounded to mean “not necessarily bounded”. We take it to mean “definitely
not bounded” instead.

– 14 –
Consider the sequence {fn }n≥1 where fn (x) := sin(2πnx). We know that sine is con-
tinuously differentiable, hence fn ∈ CC1 [0, 1] for each n ≥ 1, with

Dfn (x) = D(sin(2πnx)) = 2πn cos(2πnx).

We have
kfn k∞ = sup |fn (x)| = sup | sin(2πnx)| = sup [−1, 1] = 1
x∈[0,1] x∈[0,1]

and

kDfn k∞ = sup |Dfn (x)| = sup |2πn cos(2πnx)| = sup [−2πn, 2πn] = 2πn.
x∈[0,1] x∈[0,1]

Hence, we have
kDf k∞ kDf k∞
sup ≥ sup = sup 2πn = ∞,
f ∈C 1 [0,1] kf k∞
C
{fn }n≥1 kf k∞ n≥1

which is what we wanted. As an aside, we note that CC1 [0, 1] is not complete with respect
to the supremum norm, but it is complete with respect to the norm

kf kC 1 := kf k∞ + kf 0 k∞ .

While the derivative operator is still unbounded with respect to this new norm, in general,
the boundedness of a linear operator does depend on the choice of norms on its domain and
target, as does the numerical value of the operator norm.

Remark 2.6 . Apart from the “minor” detail that in quantum mechanics we deal with Hilbert
spaces, use a different norm than the supremum norm and that the (one-dimensional)
momentum operator acts as P (ψ) := −i~ψ 0 , the previous example is a harbinger of the fact
that the momentum operator in quantum mechanics is unbounded. This will be the case
for the position operator Q as well.

Lemma 2.7. Let (V, k·k) be a normed space. Then, addition, scalar multiplication, and the
norm are all sequentially continuous. That is, for any sequences {fn }n∈N and {gn }n∈N in
V converging to f ∈ V and g ∈ V respectively, and any sequence {zn }n∈N in C converging
to z ∈ C, we have

(i) lim (fn + gn ) = f + g

n→∞

(ii) lim zn fn = zf .
n→∞

(iii) lim kfn k = kf k

n→∞

Proof. (i) Let ε > 0. Since lim fn = f and lim gn = g by assumption, there exist
n→∞ n→∞
N1 , N2 ∈ N such that
ε
∀ n ≥ N1 : kf − fn k < 2
ε
∀ n ≥ N2 : kg − gn k < 2.

– 15 –
Let N := max{N1 , N2 }. Then, for all n ≥ N , we have

k(fn + gn ) − (f + g)k = kfn − f + gn − gk

≤ kfn − f k + kgn − gk
ε ε
< 2 + 2
= ε.

Hence lim (fn + gn ) = f + g.

n→∞

(ii) Since {zn }n∈N is a convergent sequence in C, it is bounded. That is,

∃ k > 0 : ∀ n ∈ N : |zn | ≤ k.

Let ε > 0. Since lim fn = f and lim zn = z, there exist N1 , N2 ∈ N such that
n→∞ n→∞
ε
∀ n ≥ N1 : kf − fn k <
2k
ε
∀ n ≥ N2 : kz − zn k < .
2kf k
Let N := max{N1 , N2 }. Then, for all n ≥ N , we have

kzn fn − zf k = kzn fn − zn f + zn f − zf k
= kzn (fn − f ) + (zn − z)f k
≤ kzn (fn − f )k + k(zn − z)f k
= |zn |kfn − f k + |zn − z|kf k
ε ε
<k + kf k
2k 2kf k
= ε.

Hence lim zn fn = zf .
n→∞

(iii) Let ε > 0. Since lim fn = f , there exists N ∈ N such that

n→∞

∀ n ≥ N : kfn − f k < ε.

By the triangle inequality, we have

kfn k = kfn − f + f k ≤ kfn − f k + kf k

so that kfn k − kf k ≤ kfn − f k. Similarly, kf k − kfn k ≤ kf − fn k. Since

kf − fn k = k − (fn − f )k = | − 1|kfn − f k = kfn − f k,

we have −kfn − f k ≤ kfn k − kf k ≤ kfn − f k or, by using the modulus,

kfn k − kf k ≤ kfn − f k.

Hence, for all n ≥ N , we have kfn k − kf k < ε and thus lim kfn k = kf k.

n→∞

– 16 –
Note that by taking {zn }n∈N to be the constant sequence whose terms are all equal to
some fixed z ∈ C, we have lim zfn = zf as a special case of (ii).
n→∞
This lemma will take care of some of the technicalities involved in proving the following
crucially important result.

Theorem 2.8. The set L(V, W ) of bounded linear operators from a normed space (V, k·kV )
to a Banach space (W, k · kW ), equipped with pointwise addition and scalar multiplication
and the operator norm, is a Banach space.

Proof. (a) Define addition and scalar multiplication on L(V, W ) by

(A + B)f := Af + Bf (zA)f := zAf.

It is clear that both A + B and zA are linear operators. Moreover, we have

k(A + B)f kW kAf + Bf kW
sup := sup
f ∈V kf kV f ∈V kf kV
kAf kW + kBf kW
≤ sup
f ∈V kf kV
kAf kW kBf kW
= sup + sup
f ∈V kf kV f ∈V kf kV
< ∞

since A and B are bounded. Hence, A + B is also bounded and we have

kA + Bk ≤ kAk + kBk.

Similarly, for zA we have

k(zA)f kW kzAf kW
sup := sup
f ∈V kf kV f ∈V kf kV
|z|kAf kW
= sup
f ∈V kf kV
kAf kW
= |z| sup
f ∈V kf kV
< ∞

since A is bounded and |z| is finite. Hence, zA is bounded and we have

kzAk = |z|kAk.

Thus, we have two operations

+ : L(V, W ) × L(V, W ) → L(V, W ) · : C × L(V, W ) → L(V, W )

(A, B) 7→ A + B (z, A) 7→ zA

and it is immediate to check that the vector space structure of W induces a vector
space structure on L(V, W ) with these operations.

– 17 –
(b) We need to show that (L(V, W ), k · k) is a normed space, i.e. that k · k satisfies the
properties of a norm. We have already shown two of these in part (a), namely

(b.iii) kzAk = |z|kAk

(b.iv) kA + Bk ≤ kAk + kBk.

The remaining two are easily checked.

kAf kW
(b.i) kAk := sup ≥ 0 since k · kV and k · kW are norms.
f ∈V kf kV
(b.ii) Again, by using the fat that k · kW is a norm,
kAf kW
kAk = 0 ⇔ sup =0
f ∈V kf kV
⇔ ∀ f ∈ V : kAf kW = 0
⇔ ∀ f ∈ V : Af = 0
⇔ A = 0.

Hence, (L(V, W ), k · k) is a normed space.

(c) The heart of the proof is showing that (L(V, W ), k · k) is complete. We will proceed
in three steps, analogously to the case of CC0 [0, 1].

(c.i) Let {An }n∈N be a Cauchy sequence in L(V, W ). Fix f ∈ V and let ε > 0. Then,
there exists N ∈ N such that
ε
∀ n, m ≥ N : kAn − Am k < .
kf kV
Then, for all n, m ≥ N , we have

kAn f − Am f kW = k(An − Am )f kW
k(An − Am )f kW
= kf kV
kf kV
k(An − Am )f kW
≤ kf kV sup
f ∈V kf kV
=: kf kV kAn − Am k
ε
< kf kV
kf kV
= ε.

(Note that if f = 0, we simply have kAn f − Am f kW = 0 < ε and, in the

future, we will not mention this case explicitly.) Hence, the sequence {An f }n∈N
is a Cauchy sequence in W . Since W is a Banach space, the limit limn→∞ An f
exists and is an element of W . Thus, we can define the operator

A: V → W
f 7→ lim An f,
n→∞

called the pointwise limit of {An }n∈N .

– 18 –
(c.ii) We now need to show that A ∈ L(V, W ). This is where the previous lemma
comes in handy. For linearity, let f, g ∈ V and z ∈ C. Then

A(zf + g) := lim An (zf + g)

n→∞
= lim (zAn f + An g)
n→∞
= z lim An f + lim An g
n→∞ n→∞
=: zAf + Ag

where we have used the linearity of each An and part (i) and (ii) of Lemma 2.7.
For boundedness, part (ii) and (iii) of Lemma 2.7 yield

kAf kW = lim kAn f kW

n→∞
kAn f kW
= lim kf kV
n→∞ kf kV
kAn f kW
≤ lim kf kV sup
n→∞ f ∈V kf kV
= kf kV lim kAn k
n→∞

for any f ∈ V . By rearranging, we have

kAf kW
∀f ∈ V : ≤ lim kAn k.
kf kV n→∞

Hence, to show that A is bounded, it suffices to show that the limit on the right
hand side is finite. Let ε > 0. Since {An }n∈N is a Cauchy sequence, there exists
N ∈ N such that
∀ n, m ≥ N : kAn − Am k < ε.
Then, by the proof of part (i) of Lemma 2.7, we have

kAn k − kAm k ≤ kAn − Am k < ε

for all n, m ≥ N . Hence, the sequence of real numbers {kAn k}n∈N is a Cauchy
sequence. Since R is complete, this sequence converges to some real number
r ∈ R. Therefore
kAf kW
sup ≤ lim kAn k = r < ∞
f ∈V kf kV n→∞

and thus A ∈ L(V, W ).

(c.iii) To conclude, we have to show that lim An = A. Let ε > 0. Then
n→∞

kAn − Ak = kAn + Am − Am − Ak ≤ kAn − Am k + kAm − Ak.

Since {An }n∈N is Cauchy, there exists N1 ∈ N such that

∀ n, m ≥ N1 : kAn − Am k < 2ε .

– 19 –
Moreover, since A is the pointwise limit of {An }n∈N , for any f ∈ V there exists
N2 ∈ N such that
εkf kV
∀ m ≥ N2 : kAm f − Af kW <
2
and hence, for all m ≥ N2
εkf k
V
kAm f − Af kW ε
kAm − Ak := sup ≤ 2 =
f ∈V kf kV kf kV 2

Let N := max{N1 , N2 } and fix m ≥ N . Then, for all n ≥ N , we have

ε ε
kAn − Ak ≤ kAn − Am k + kAm − Ak < 2 + 2 = ε.

Thus, lim An = A and we call A the uniform limit of {An }n∈N .

n→∞
This concludes the proof that (L(V, W ), k · k) is a Banach space.

Remark 2.9 . Note that if V and W are normed spaces, then L(V, W ) is again a normed
space, while for L(V, W ) to be a Banach space it suffices that W be a Banach space.

Remark 2.10 . In the proof that L(V, W ) is a Banach space, we have shown a useful in-
equality which we restate here for further emphasis. If A : V → W is bounded, then

∀ f ∈ V : kAf kW ≤ kAkkf kV

The following is an extremely important special case of L(V, W ).

Definition. Let V be a normed space. Then V ∗ := L(V, C) is called the dual of V .

Note that, since C is a Banach space, the dual of a normed space is a Banach space.
The elements of V ∗ are variously called covectors or functionals on V .

Remark 2.11 . You may recall from undergraduate linear algebra that the dual of a vector
space was defined to be the vector space of all linear maps V → C, rather than just the
bounded ones. This is because, in finite dimensions, all linear maps are bounded. So the
two definitions agree as long as we are in finite dimensions. If we used the same definition
for the infinite-dimensional case, then V ∗ would lack some very desirable properties, such
as that of being a Banach space.

The dual space can be used to define a weaker notion of convergence called, rather
unimaginatively, weak convergence.

Definition. A sequence {fn }n∈N is said to converge weakly to f ∈ V if

∀ ϕ ∈ V ∗ : lim ϕ(fn ) = ϕ(f ).

n→∞

– 20 –
Note that {ϕ(fn )}n∈N is just a sequence of complex numbers. To indicate that the
sequence {fn }n∈N converges weakly to f ∈ V we write

w-lim fn = f.
n→∞

In order to further emphasise the distinction with weak convergence, we may say that
{fn }n∈N converges strongly to f ∈ V if it converges according to the usual definition, and
we will write accordingly
s-lim fn = f.
n→∞

Proposition 2.12. Let {fn }n∈N be a sequence in a normed space (V, k · kV ). If {fn }n∈N
converges strongly to f ∈ V , then it also converges weakly to f ∈ V , i.e.

s-lim fn = f ⇒ w-lim fn = f.
n→∞ n→∞

Proof. Let ε > 0 and let ϕ ∈ V ∗ . Since {fn }n∈N converges strongly to f ∈ V , there exists
N ∈ N such that
ε
∀ n ≥ N : kfn − f kV < .
kϕk
Then, since ϕ ∈ V ∗ is bounded, we have

|ϕ(fn ) − ϕ(f )| = |ϕ(fn − f )|

≤ kϕkkfn − f kV
ε
< kϕk
kϕk
=ε

for any n ≥ N . Hence, lim ϕ(fn ) = ϕ(f ). That is, w-lim fn = f .

n→∞ n→∞

2.3 Extension of bounded linear operators

Note that, so far, we have only considered bounded linear maps A : DA → W where DA is
the whole of V , rather than a subspace thereof. The reason for this is that we will only
consider densely defined linear maps in general, and any bounded linear map from a dense
subspace of V can be extended to a bounded linear map from the whole of V . Moreover,
the extension is unique. This is the content of the so-called BLT6 theorem.

Lemma 2.13. Let (V, k · k) be a normed space and let DA be a dense subspace of V . Then,
for any f ∈ V , there exists a sequence {αn }n∈N in DA which converges to f .

Proof. Let f ∈ V . Clearly, there exists a sequence {fn }n∈N in V which converges to f (for
instance, the constant sequence). Let ε > 0. Then, there exists N ∈ N such that

∀ n ≥ N : kfn − f k < 2ε .
6
Bounded Linear Transformation, not Bacon, Lettuce, Tomato.

– 21 –
Since DA is dense in V and each fn ∈ V , we have
∀ n ∈ N : ∃ αn ∈ DA : kαn − fn k < 2ε .

The sequence {αn }n∈N is a sequence in DA and we have

kαn − f k = kαn − fn + fn − f k
≤ kαn − fn k + kfn − f k
ε ε
< 2 + 2
=ε
for all n ≥ N . Hence lim αn = f .
n→∞

Definition. Let V, W be vector spaces and let A : DA → W be a linear map, where DA ⊆ V .

An extension of A is a linear map A
b : V → W such that

∀ α ∈ DA : Aα
b = Aα.

Theorem 2.14 (BLT theorem). Let V be a normed space and W a Banach space. Any
b : V → W such that A
densely defined linear map A : DA → W has a unique extension A b is
bounded. Moreover, kAk b = kAk.

Proof. (a) Let A ∈ L(DA , W ). Since DA is dense in V , for any f ∈ V there exists a
sequence {αn }n∈N in DA which converges to f . Moreover, since A is bounded, we
have
∀ n ∈ N : kAαn − Aαm kW ≤ kAkkαn − αm kV ,
from which it quickly follows that {Aαn }n∈N is Cauchy in W . As W is a Banach
space, this sequence converges to an element of W and thus we can define
b: V → W
A
f 7→ lim Aαn ,
n→∞
where {αn }n∈N is any sequence in DA which converges to f .

(b) First, let us show that A

b is well-defined. Let {αn }n∈N and {βn }n∈N be two sequences
in DA which converge to f ∈ V and let ε > 0. Then, there exist N1 , N2 ∈ N such that
ε
∀ n ≥ N1 : kαn − f kV <
2kAk
ε
∀ n ≥ N2 : kβn − f kV < .
2kAk
Let N := max{N1 , N2 }. Then, for all n ≥ N , we have
kAαn − Aβn kW = kA(αn − βn )kW
≤ kAkkαn − βn kV
= kAkkαn − f + f − βn kV
≤ kAk(kαn − f kV + kf − βn kV )
ε ε
< kAk +
2kAk 2kAk
= ε,

– 22 –
where we have used the fact that A is bounded. Thus, we have shown

lim (Aαn − Aβn ) = 0.

n→∞

Then, by using Lemma 2.7 and rearranging, we find

lim Aαn = lim Aβn ,

n→∞ n→∞

that is, A
b is indeed well-defined.

(c) To see that A

b is an extension of A, let α ∈ DA . The constant sequence {αn }n∈N with
αn = α for all n ∈ N is a sequence in DA converging to α. Hence

Aα
b := lim Aαn = lim Aα = Aα.
n→∞ n→∞

(d) We now check that A ∈ L(V, W ). For linearity, let f, g ∈ V and z ∈ C. As DA is

dense in V , there exist sequences {αn }n∈N and {βn }n∈N in DA converging to f and
g, respectively. Moreover, as DA is a subspace of V , the sequence {γn }n∈N given by

γn := zαn + βn

is again a sequence in DA and, by Lemma 2.7,

lim γn = zf + g.
n→∞

Then, we have

A(zf
b + g) := lim Aγn
n→∞
= lim A(zαn + βn )
n→∞
= lim (zAαn + Aβn )
n→∞
= z lim Aαn + lim Aβn
n→∞ n→∞
=: z Af
b + Ag.
b

For boundedness, let f ∈ V and {αn }n∈N a sequence in DA which converges to f .

Then, since A is bounded,

kAf
b kW := lim Aαn
W
n→∞
= lim kAαn kW
n→∞
≤ lim kAkkαn kV
n→∞
= kAk lim kαn kV
n→∞
= kAkkf kV .

Therefore
kAf
b kW kAkkf kV
sup ≤ sup = sup kAk = kAk < ∞
f ∈V kf kV f ∈V kf kV f ∈V

and hence A
b is bounded.

– 23 –
(e) For uniqueness, suppose that A
e ∈ L(V, W ) is another extension of A. Let f ∈ V and
{αn }n∈N a sequence in DA which converges to f . Then, we have

kAf
e − Aαn kW = kAf
e − Aα
e n kW ≤ kAkkf
e − αn kV .

It follows that
e − Aαn ) = 0
lim (Af
n→∞

and hence, for all f ∈ V ,

Af
e = lim Aαn =: Af.
b
n→∞

Therefore, A
e = A.
b

(f) Finally, we have already shown in part (d) that

kAf
b kW
kAk
b := sup ≤ kAk.
f ∈V kf kV

On the other hand, since DA ⊆ V , we must also have

kAf kW kAf
b kW kAf
b kW
kAk := sup = sup ≤ sup =: kAk.
b
f ∈DA kf kV f ∈DA kf kV f ∈V kf kV

Hence, we also have kAk ≤ kAk.

b Thus, kAk
b = kAk.

Remark 2.15 . Note a slight abuse of notation in the equality kAk

b = kAk. The linear maps
A and A belong to L(V, W ) and L(DA , W ), respectively. These are different normed (in
b
fact, Banach) spaces and, in particular, carry different norms. To be more precise, we
should have written
kAk
b L(V,W ) = kAkL(D ,W ) ,
A

where
kAf
b kW kAf kW
kAk
b L(V,W ) := sup and kAkL(DA ,W ) := sup .
f ∈V kf kV f ∈DA kf kV

– 24 –
3 Separable Hilbert Spaces

3.1 Relationship between norms and inner products

A Hilbert space is a vector space (H, +, ·) equipped with a sesqui-linear inner product h·|·i
which induces a norm k · kH with respect to which H is a Banach space. Note that by “being
induced by h·|·i” we specifically mean that the norm is defined as

k · k: V → R
p
f 7→ hf |f i.

Recall that a sesqui-linear inner product on H is a map h·|·i : H × H → C which is con-

jugate symmetric, linear in the second argument and positive-definite. Note that conjugate
symmetry together with linearity in the second argument imply conjugate linearity in the
first argument:

hzψ1 + ψ2 |ϕi = hϕ|zψ1 + ψ2 i

Of course, since Hilbert spaces are a special case of Banach spaces, everything that we
have learned about Banach spaces also applies to Hilbert paces. For instance, L(H, H),
the collection of all bounded linear maps H → H, is a Banach space with respect to the
operator norm. In particular, the dual of a Hilbert space H is just H∗ := L(H, C). We will
see that the operator norm on H∗ is such that there exists an inner product on H∗ which
induces it, so that the dual of a Hilbert space is again a Hilbert space.
First, in order to check that the norm induced by an inner product on V is indeed a
norm on V , we need one of the most important inequalities in mathematics.

Proposition 3.1 (Cauchy-Schawrz inequality7 ). Let h·|·i be a sesqui-linear inner product

on V . Then, for any f, g ∈ V , we have

|hf |gi|2 ≤ hf |f ihg|gi.

Proof. If f = 0 or g = 0, then equality holds. Hence suppose that f 6= 0 and let

hf |gi
z := ∈ C.
hf |f i
7
Also known as the Cauchy-Bunyakovsky-Schwarz inequality in the Russian literature.

– 25 –
Then, by positive-definiteness of h·|·i, we have

0 ≤ hzf − g|zf − gi
= |z|2 hf |f i − zhf |gi − zhg|f i + hg|gi
|hf |gi|2 hf |gi hf |gi
= 2
hf |f i − hf |gi − hf |gi + hg|gi
hf |f i hf |f i hf |f i
|hf |gi|2 |hf |gi|2 |hf |gi|2
= − − + hg|gi
hf |f i hf |f i hf |f i
|hf |gi|2
=− + hg|gi.
hf |f i
By rearranging, since hf |f i > 0, we obtain the desired inequality.

Note that, by defining kf k := hf |f i, and using the fact that |hf |gi| ≥ 0, we can write
p

the Cauchy-Schwarz inequality as

|hf |gi| ≤ kf kkgk.

Proposition 3.2. The induced norm on V is a norm.

Proof. Let f, g ∈ V and z ∈ C. Then

(i) kf k := hf |f i ≥ 0
p

(ii) kf k = 0 ⇔ kf k2 = 0 ⇔ hf |f i = 0 ⇔ f = 0 by positive-definiteness

(iii) kzf k := hzf |zf i = zzhf |f i = |z|2 hf |f i = |z| hf |f i =: |z|kf k

p p p p

(iv) Using the fact that z+z = 2 Re z and Re z ≤ |z| for any z ∈ C and the Cauchy-Schwarz
inequality, we have

By taking the square root of both sides, we have kf + gk ≤ kf k + kgk.

Hence, we see that any inner product space (i.e. a vector space equipped with a sesqui-
linear inner product) is automatically a normed space under the induced norm. It is only
natural to wonder whether the converse also holds, that is, whether every norm is induced
by some sesqui-linear inner product. Unfortunately, the answer is negative in general. The
following theorem gives a necessary and sufficient condition for a norm to be induced by a
sesqui-linear inner product and, in fact, by a unique such.

– 26 –
Theorem 3.3 (Jordan-von Neumann). Let V be a vector space. A norm k · k on V is
induced by a sesqui-linear inner product h·|·i on V if, and only if, the parallelogram identity

kf + gk2 + kf − gk2 = 2kf k2 + 2kgk2

holds for all f, g ∈ V , in which case, h·|·i is determined by the polarisation identity
3
1X k
hf |gi = i kf + i4−k gk2
4
k=0
1
= (kf + gk2 − kf − gk2 + ikf − igk2 − ikf + igk2 ).
4
Proof. (⇒) If k · k is induced by h·|·i, then by direct computation

kf + gk2 + kf − gk2 := hf + g|f + gi + hf − g|f − gi

so the parallelogram identity is satisfied. We also have

kf + gk2 − kf − gk2 := hf + g|f + gi − hf − g|f − gi

and

Therefore
kf + gk2 − kf − gk2 + ikf − igk2 − ikf + igk2 = 4hf |gi.
that is, the inner product is determined by the polarisation identity.

(⇐) Suppose that k · k satisfies the parallelogram identity. Define h·|·i by

1
hf |gi := (kf + gk2 − kf − gk2 + ikf − igk2 − ikf + igk2 ).
4
We need to check that this satisfies the defining properties of a sesqui-linear inner
product.

– 27 –
(i) For conjugate symmetry
1

hf |gi := 2 2 2 2
4 kf + gk − kf − gk + ikf − igk − ikf + igk
1 2 2 2 2
= 4 (kf + gk − kf − gk − ikf − igk + ikf + igk )
1 2 2 2 2
= 4 (kf + gk − kf − gk − ik(−i)(if + g)k + iki(−if + g)k )
1 2 2 2 2 2 2
= 4 (kg + f k − kg − f k − i(| − i|) kg + if k + i(|i|) kg − if k )
1 2 2 2 2
= 4 (kg + f k − kg − f k − ikg + if k + ikg − if k )
=: hg|f i

(ii) We will now show linearity in the second argument. This is fairly non-trivial
and quite lengthy. We will focus on additivity first. We have
1
hf |g + hi := (kf + g + hk2 − kf − g − hk2 + ikf − ig − ihk2 − ikf + ig + ihk2 ).
4
Consider the real part of hf |g + hi. By successive applications of the parallelo-
gram identity, we find

Rehf |g + hi = 14 (kf + g + hk2 − kf − g − hk2 )

= 41 (kf + g + hk2 + kf + g − hk2 − kf + g − hk2 − kf − g − hk2 )
= 41 (2kf + gk2 + 2khk2 − 2kf − hk2 − 2kgk2 )
= 14 (2kf + gk2 + 2kf k2 + 2khk2 − 2kf − hk2 − 2kf k2 − 2kgk2 )
= 14 (2kf + gk2 + kf + hk2 + kf − hk2 − 2kf − hk2 − kf + gk2 − kf − gk2 )
= 41 (kf + gk2 + kf + hk2 − kf − hk2 − kf − gk2 )
= Rehf |gi + Rehf |hi.

Replacing g and h with −ig and −ih respectively, we obtain

Imhf |g + hi = Imhf |gi + Imhf |hi.

Hence, we have

which proves additivity.

For scaling invariance, we will proceed in several steps.
(a) First, note that

hf |0i := 14 (kf k2 − kf k2 + ikf k2 − ikf k2 ) = 0

and hence hf |0gi = 0hf |gi holds.

– 28 –
(b) Suppose that hf |ngi = nhf |gi for some n ∈ N. Then, by additivity

hf |(n + 1)gi = hf |ng + gi

Hence, by induction on n with base case (a), we have

∀ n ∈ N : hf |ngi = nhf |gi.

(c) Note that by additivity

(a)
hf |gi + hf |−gi = hf |g − gi = hf |0i = 0.

Hence hf |−gi = −hf |gi.

(d) Then, for any n ∈ N
(c) (b)
hf |−ngi = −hf |ngi = −nhf |gi

and thus
∀ n ∈ Z : hf |ngi = nhf |gi.
(e) Now note that for any m ∈ Z \ {0}

1 (d) 1
mhf | m gi = hf |m m gi = hf |gi

and hence, by dividing by m, we have hf | m

1
gi = 1
m hf |gi.
(f) Therefore, for any r = n
m ∈ Q, we have

Rehf |gi = 14 (kf + gk2 − kf − gk2 )

= 14 (2kf + gk2 − kf + gk2 − kf − gk2 )
= 41 (2kf + gk2 − 2kf k2 − 2kgk2 )
≤ 14 (2(kf k + kgk)2 − 2kf k2 − 2kgk2 )
= 14 (2kf k2 + 4kf kkgk + 2kgk2 − 2kf k2 − 2kgk2 )
= kf kkgk.

– 29 –
Replacing g with −ig and noting that k − igk = | − i|kgk = kgk, we also have

Imhf |gi ≤ kf kkgk.

Hence, we find

|hf |gi| = | Rehf |gi + i Imhf |gi|

p
= (Rehf |gi)2 + (Imhf |gi)2
p
≤ (kf kkgk)2 + (kf kkgk)2
√
= 2kf kkgk.

(h) Let r ∈ R. Since R is the completion of Q (equivalently, Q is dense in R),

there exists a sequence {rn }n∈N in Q which converges to r. Let ε > 0. Then,
there exist N1 , N2 ∈ N such that
ε
∀ n ≥ N1 : |rn − r| < √
2 2kf kkgk
ε
∀ n, m ≥ N2 : |rn − rm | < √ .
2 2kf kkgk
Let N := max{N1 , N2 } and fix m ≥ N . Then, for all n ≥ N , we have

|rn hf |gi − hf |rgi| = |rn hf |gi − rm hf |gi + rm hf |gi − hf |rgi|

that is, lim rn hf |gi = hf |rgi.

n→∞
(i) Hence, for any r ∈ R, we have
(h)
rhf |gi = lim rn hf |gi = lim rn hf |gi = hf |rgi
n→∞ n→∞

and thus
∀ r ∈ R : rhf |gi = hf |rgi.
(j) We now note that

hf |igi := 14 (kf + igk2 − kf − igk2 + ikf − i2 gk2 − ikf + i2 gk2 )

= 1
4 i (−ikf + igk2 + kf − igk2 + ikf + gk2 − kf − gk2 )
=: ihf |gi

and hence hf |igi = ihf |gi.

Hence, h·|·i is indeed a sesqui-linear inner product. Note that, from part (iii) above,
we have p
hf |f i = kf k.
That is, the inner product h·|·i does induce the norm from which we started, and this
completes the proof.

Remark 3.4 . Our proof of linearity is based on the hints given in Section 6.1, Exercise
27, from Linear Algebra (4th Edition) by Friedberg, Insel, Spence. Other proofs of the
Jordan-von Neumann theorem can be found in
• Kadison, Ringrose, Fundamentals of the Theory of Operator Algebras: Volume I:
Elementary Theory, American Mathematical Society 1997

• Kutateladze, Fundamentals of Functional Analysis, Springer 1996.

• Day, Normed Linear Spaces, Springer 1973.

Remark 3.5 . Note that, often in the more mathematical literature, a sesqui-linear inner
product is defined to be linear in the first argument rather than the second. In that case,
the polarisation identity takes the form
3
1X k
hf |gi = i kf + ik gk2
4
k=0
1
= (kf + gk2 − kf − gk2 + ikf + igk2 − ikf − igk2 ).
4

– 31 –
Example 3.6 . Consider CC0 [0, 1] and let f (x) = x and g(x) = 1. Then

kf k∞ = 1, kgk∞ = 1, kf + gk∞ = 2, kf − gk∞ = 1

and hence
kf + gk2∞ + kf − gk2∞ = 5 6= 4 = 2kf k2∞ + 2kgk2∞ .
Thus, by the Jordan-von Neumann theorem, there is no inner product on CC0 [0, 1] which
induces the supremum norm. Therefore, (CC0 [0, 1], k · k∞ ) cannot be a Hilbert space.

Proposition 3.7. Let H be a Hilbert space. Then, H∗ is a Hilbert space.

Proof. We already know that H∗ := L(H, C) is a Banach space. The norm on H∗ is just
the usual operator norm
|f (ϕ)|
kf kH∗ := sup
ϕ∈H kϕkH

where, admittedly somewhat perversely, we have reversed our previous notation for the dual
√
elements. Since the modulus is induced by the standard inner product on C, i.e. |z| = zz,
it satisfies the parallelogram identity. Hence, we have

|(f + g)(ϕ)| 2 |(f − g)(ϕ)| 2

2 2
kf + gkH∗ + kf − gkH∗ := sup + sup
ϕ∈H kϕkH ϕ∈H kϕkH
|(f + g)(ϕ)|2 |(f − g)(ϕ)|2
= sup + sup
ϕ∈H kϕk2H ϕ∈H kϕk2H
|f (ϕ) + g(ϕ)|2 + |f (ϕ) − g(ϕ)|2
= sup
ϕ∈H kϕk2H
2|f (ϕ)|2 + 2|g(ϕ)|2
= sup
ϕ∈H kϕk2H
|f (ϕ)|2 |g(ϕ)|2
= 2 sup 2 + 2 sup 2
ϕ∈H kϕkH ϕ∈H kϕkH
=: 2kf k2H∗ + 2kgk2H∗ ,

where several steps are justified by the fact that the quantities involved are non-negative.
Hence, by the Jordan-von Neumann theorem, the inner product on H∗ defined by the
polarisation identity induces k · kH∗ . Hence, H∗ is a Hilbert space.

The following useful fact is an immediate application of the Cauchy-Schwarz inequality.

Proposition 3.8. Inner products on a vector space are sequentially continuous.

Proof. Let h·|·i be an inner product on V . Fix ϕ ∈ V and let lim ψn = ψ. Then
n→∞

|hϕ|ψn i − hϕ|ψi| = |hϕ|ψn − ψi|

≤ kϕkkψn − ψk

and hence lim hϕ|ψn i = hϕ|ψi.

n→∞

– 32 –
3.2 Hamel versus Schauder8
Choosing a basis on a vector space is normally regarded as mathematically inelegant. The
reason for this is that most statements about vector spaces are much clearer and, we main-
tain, aesthetically pleasing when expressed without making reference to a basis. However,
in addition to the fact that some statements are more easily and usefully written in terms
of a basis, bases provide a convenient way to specify the elements of a vector space in terms
of components. The notion of basis for a vector space that you most probably met in your
linear algebra course is more properly know as Hamel basis.

Definition. A Hamel basis of a vector space V is a subset B ⊆ V such that

(i) any finite subset {e1 , . . . , en } ⊆ B is linearly independent, i.e.

n
X
λi ei = 0 ⇒ λ1 = · · · = λn = 0
i=1

(ii) the set B is a generating (or spanning) set for V . That is, for any element v ∈ V ,
there exist a finite subset {e1 , . . . , en } ⊆ B and λ1 , . . . , λn ∈ C such that
n
X
v= λi ei .
i=1

Equivalently, by defining the linear span of a subset U ⊆ V as

X n
λ ui λ , . . . , λ ∈ C, u1 , . . . , un ∈ U and n ≥ 1 ,
i 1 n
span U :=
i=1

i.e. the set of all finite linear combinations of elements of U with complex coefficients,
we can restate this condition simply as V = span B.

Given a basis B, one can show that for each v ∈ V the λ1 , . . . , λn appearing in (ii)
above are uniquely determined. They are called the components of v with respect to B.
One can also show that if a vector space admits a finite Hamel basis B, then any other
basis of V is also finite and, in fact, of the same cardinality as B.

Definition. If a vector space V admits a finite Hamel basis, then it is said to be finite-
dimensional and its dimension is dim V := |B|. Otherwise, it is said to be infinite-
dimensional and we write dim V = ∞.

Theorem 3.9. Every vector space admits a Hamel basis.

For a proof of (a slightly more general version of) this theorem, we refer the interested
reader to Dr Schuller’s Lectures on the Geometric Anatomy of Theoretical Physics.
Note that the proof that every vector space admits a Hamel basis relies on the axiom
of choice and, hence, it is non-constructive. By a corollary to Baire’s category theorem, a
8
Not a boxing competition.

– 33 –
Hamel basis on a Banach space is either finite or uncountably infinite. Thus, while every
Banach space admits a Hamel basis, such bases on infinite-dimensional Banach spaces are
difficult to construct explicitly and, hence, not terribly useful to express vectors in terms
of components and perform computations. Thankfully, we can use the extra structure of a
Banach space to define a more useful type of basis.

Definition. Let (W, k · k) be a Banach space. A Schauder basis of W is a sequence {en }n∈N
in W such that, for any f ∈ W , there exists a unique sequence {λn }n∈N in C such that
n
X ∞
X
f = lim λi ei =: λi ei
n→∞
i=0 i=0

or, by explicitly using the definition of limit in W ,

n
X

i

lim f −
λ ei = 0.
n→∞
i=0

We note the following points.

• Since Schauder bases require a notion of convergence, they can only be defined on
a vector space equipped with a (compatible) topological structure, of which Banach
spaces are a special case.

• Unlike Hamel bases, Schauder bases need not exist.

• Since the convergence of a series may depend on the order of its terms, Schauder bases
must be considered as ordered bases. Hence, two Schauder bases that merely differ
in the ordering of their elements are different bases, and permuting the elements of a
Schauder basis doesn’t necessarily yield another Schauder basis.

• The uniqueness requirement in the definition immediately implies that the zero vector
cannot be an element of a Schauder basis.

• Schauder bases satisfy a stronger linear independence property than Hamel bases,
namely
X∞
λi ei = 0 ⇒ ∀ i ∈ N : λi = 0.
i=0

• At the same time, they satisfy a weaker spanning condition. Rather than the linear
span of the basis being equal to W , we only have that it is dense in W . Equivalently,

W = span{en | n ∈ N},

where the topological closure U of a subset U ⊆ W is defined as

U := lim un | ∀ n ∈ N : un ∈ U .
n→∞

– 34 –
Definition. A Schauder basis {en }n∈N of (W, k · k) is said to be normalised if

∀ n ∈ N : ken k = 1.

Multiplying an element of a Schauder basis by a complex number gives again a Schauder

basis (not the same one, of course). Since Schauder bases do not contain the zero vector,
any Schauder basis {en }n∈N gives rise to a normalised Schauder basis {e
en }n∈N by defining
en
een := .
ken k

3.3 Separable Hilbert spaces

Separability is a topological property. A topological space is said to be separable if it
contains a dense subset which is also countable. A Banach space is said to be separable
if it is separable as a topological space with the topology induced by the norm. Similarly,
a Hilbert space is said to be separable if it is separable as a topological space with the
topology induced by the norm induced in turn by the inner product.
For infinite-dimensional Hilbert spaces, there is a much more useful characterisation of
separability, which we will henceforth take as our definition.

Proposition 3.10. An infinite-dimensional Hilbert space is separable if, and only if, it
admits an orthonormal Schauder basis. That is, a Schauder basis {en }n∈N such that
(
1 if i = j
∀ i, j ∈ N : hei |ej i = δij := .
0 if i 6= j

Whether this holds for Banach spaces or not was a famous open problem in
functional analysis, problem 153 from the Scottish book. It was solved in 1972,
more that three decades after it was first posed, when Swedish mathematician
Enflo constructed an infinite-dimensional separable Banach space which lacks
a Schauder basis. That same year, he was awarded a live goose9 for his effort.

Remark 3.11 . The Kronecker symbol δij appearing above does not represent the compo-
nents of the identity map on H. Instead, δij are the components of the sesqui-linear form
h·|·i, which is a map H × H → C, unlike idH which is a map H → H. If not immediately
understood, this remark may be safely ignored.

Remark 3.12 . In finite-dimensions, since every vector space admits (by definition) a finite
Hamel basis, every inner product space admits an orthonormal basis by the Gram-Schmidt
orthonormalisation process.

From now on, we will only consider orthonormal Schauder bases, sometimes also called
Hilbert bases, and just call them bases.

Lemma 3.13. Let H be a Hilbert space with basis {en }n∈N . The unique sequence in the
expansion of ψ ∈ H in terms of this basis is {hen |ψi}n∈N .
9
https://en.wikipedia.org/wiki/Per_Enflo#Basis_problem_of_Banach

– 35 –
Proof. By using the continuity of the inner product, we have
∞ j
X
hei |ψi = ei
λ ej
j=0
n
X
j

= ei lim λ ej
n→∞
j=0
n j
X
= lim ei λ ej
n→∞
j=0
n
X
= lim λj hei |ej i
n→∞
i=0
∞
X
= λj δij
i=0
= λi ,

which is what we wanted.

While we have already used the term orthonormal, let us note that this means both
orthogonal and normalised. Two vectors ϕ, ψ ∈ H are said to be orthogonal if hϕ|ψi = 0,
and a subset of H is called orthogonal if its elements are pairwise orthogonal.

Lemma 3.14 (Pythagoras’ theorem). Let H be a Hilbert space and let {ψ0 , . . . , ψn } ⊂ H
be a finite orthogonal set. Then

n n
X 2 X
kψi k2 .

ψi =

i=0 i=0

Proof. Using the pairwise orthogonality of {ψ0 , . . . , ψn }, we simply calculate

n n
n n n
X 2 X X X X X
kψi k2 .

ψi
:= ψi
ψj = hψ |ψ
i i i + hψ |ψ
i j i =:
i=0 i=0 j=0 i=0 i6=j i=0

Corollary 3.15. Let ψ ∈ H and let {en }n∈N be a basis of H. Then

∞
X
kψk2 = |hei |ψi|2 .
i=0

Proof. By continuity of the norm, we have

n 2 n ∞
2
X X
2 2
X
|hei |ψi|2 .

kψk = lim hei |ψiei = lim
|hei |ψi| kei k =
n→∞ n→∞
i=0 i=0 i=0

3.4 Unitary maps

An insightful way to study a structure in mathematics is to consider maps between different
instances A, B, C, . . . of that structure, and especially the structure-preserving maps. If

– 36 –
a certain structure-preserving map A → B is invertible and its inverse is also structure-
preserving, then both these maps are generically called isomorphisms and A and B are
said to be isomorphic instances of that structure. Isomorphic instances of a structure are
essentially the same instance of that structure, just dressed up in different ways. Typically,
there are infinitely many concrete instances of any given structure. The highest form of
understanding of a structure that we can hope to achieve is that of a classification of its
instances up to isomorphism. That is, we would like to know how many different, non-
isomorphic instances of a given structure there are.
In linear algebra, the structure of interest is that of vector space over some field F. The
structure-preserving maps are just the linear maps and the isomorphisms are the linear
bijections (whose inverses are automatically linear). Finite-dimensional vector spaces over
F are completely classified by their dimension, i.e. there is essentially only one vector space
over F for each n ∈ N, and Fn is everyone’s favourite. Assuming the axiom of choice,
infinite-dimensional vector spaces over F are classified in the same way, namely, there is,
up to linear isomorphism, only one vector space over F for each infinite cardinal.
Of course, one could do better and also classify the base fields themselves. The classi-
fication of finite fields (i.e. fields with a finite number of elements) was achieved in 1893 by
Moore, who proved that the order (i.e. cardinality) of a finite field is necessarily a power of
some prime number, and there is only one finite field of each order, up to the appropriate
notion of isomorphism. The classification of infinite fields remains an open problem.
A classification with far-reaching implications in physics is that of finite-dimensional,
semi-simple, complex Lie algebras, which is discussed in some detail in Dr Schuller’s Lectures
on the Geometric Anatomy of Theoretical Physics.
The structure-preserving maps between Hilbert spaces are those that preserve both the
vector space structure and the inner product. The Hilbert space isomorphisms are called
unitary maps.
Definition. Let H and G be Hilbert spaces. A bounded bijection U ∈ L(H, G) is called a
unitary map (or unitary operator ) if
∀ ψ, ϕ ∈ H : hU ψ|U ϕiG = hψ|ϕiH .
If there exists a unitary map H → G, then H and G are said to be unitarily equivalent and
we write H ∼ =Hil G.
There are a number of equivalent definitions of unitary maps (we will later see one
involving adjoints) and, in fact, our definition is fairly redundant.
Proposition 3.16. Let U : H → G be a surjective map which preserves the inner product.
Then, U is a unitary map.
Proof. (i) First, let us check that U is linear. Let ψ, ϕ ∈ H and z ∈ C. Then
kU (zψ + ϕ) − zU ψ − U ϕk2G = hU (zψ + ϕ) − zU ψ − U ϕ|U (zψ + ϕ) − zU ψ − U ϕiG
= hU (zψ + ϕ)|U (zψ + ϕ)iG + |z|2 hU ψ|U ψiG + hU ϕ|U ϕiG
− zhU (zψ + ϕ)|U ψiG − hU (zψ + ϕ)|U ϕiG + zhU ψ|U ϕiG
− zhU ψ|U (zψ + ϕ)iG − hU ϕ|U (zψ + ϕ)iG + zhU ϕ|U ψiG

Hence kU (zψ + ϕ) − zU ψ − U ϕkG = 0, and thus

U (zψ + ϕ) = zU ψ + U ϕ.

(ii) For boundedness, simply note that since for any ψ ∈ H

p p
kU ψkG := hU ψ|U ψiG = hψ|ψiH =: kψkH ,

we have
kU ψkG
sup = 1 < ∞.
ψ∈H kψkH

Hence U is bounded and, in fact, has unit operator norm.

(iii) Finally, recall that a linear map is injective if, and only if, its kernel is trivial. Suppose
that ψ ∈ ker U . Then, we have

hψ|ψiH = hU ψ|U ψiG = h0|0iG = 0.

Hence, by positive-definiteness, ψ = 0 and thus, U is injective. Since U is also

surjective by assumption, it satisfies our definition of unitary map.

Note that a map U : H → G is called an isometry if

∀ ψ ∈ H : kU ψkG = kψkH .

Linear isometries are, of course, the structure-preserving maps between normed spaces. We
have shown that every unitary map is an isometry has unit operator norm, hence the name
unitary operator.

Example 3.17 . Consider the set of all square-summable complex sequences

∞
X
2 2
` (N) := a : N → C |ai | < ∞ .

i=0

We define addition and scalar multiplication of sequences termwise, that is, for all n ∈ N
and all complex numbers z ∈ C,

(a + b)n := an + bn
(z · a)n := zan .

– 38 –
These are, of course, just the discrete analogues of pointwise addition and scalar multipli-
cation of maps. The triangle inequality and homogeneity of the modulus, together with the
vector space structure of C, imply that (`2 (N), +, ·) is a complex vector space.
The standard inner product on `2 (N) is
∞
X
ha|bi`2 := ai bi .
i=0

This inner product induces the norm

v
u∞
p uX
kak`2 := ha|ai`2 =t |ai |2 ,
i=0

with respect to which `2 (N) is complete. Hence, (`2 (N), +, ·, h·|·i`2 ) is a Hilbert space.
Consider the sequence of sequences {en }n∈N where

e0 = (1, 0, 0, 0, . . .)
e1 = (0, 1, 0, 0, . . .)
e2 = (0, 0, 1, 0, . . .)
..
.

i.e. we have (en )m = δnm . Each a ∈ `2 (N) can be written uniquely as

∞
X
a= λi ei ,
i=0

where λi = hei |ai`2 = ai . The sequences en are clearly square-summable and, in fact, they
are orthonormal with respect to h·|·i`2
∞
X ∞
X
hen |em i`2 := (en )i (em )i = δni δmi = δnm .
i=0 i=0

Hence, the sequence {en }n∈N is an orthonormal Schauder basis of `2 (N), which is
therefore an infinite-dimensional separable Hilbert space.

Theorem 3.18 (Classification of separable Hilbert spaces). Every infinite-dimensional sep-

arable Hilbert space is unitarily equivalent to `2 (N).

Proof. Let H be a separable Hilbert space with basis {en }n∈N . Consider the map

U : H → `2 (N)
ψ 7→ {hen |ψiH }n∈N .

Note that, for any ψ ∈ H, the sequence {hen |ψiH }n∈N is indeed square-summable since we
have
∞
X
|hei |ψiH |2 = kψk2H < ∞.
i=0

– 39 –
By our previous proposition, in order to show that U is a unitary map, it suffices to
show that it is surjective and preserves the inner product. For surjectivity, let {an }n∈N
be a complex square-summable sequence. Then, by elementary analysis, we know that
lim |an | = 0. This implies that, for any ε > 0, there exists N ∈ N such that
n→∞

n
X
∀n ≥ m ≥ N : |ai |2 < ε.
i=m

Then, for all n, m ≥ N (without loss of generality, assume n > m), we have
n m 2 n 2 n n
X X X X
2 2
X
|ai |2 < ε.

a e
i i − a e
j j
=
a e
i i
= |ai | ke k
i H =
i=0 j=0 H i=m+1 H i=m+1 i=m+1
Pn
That is, i=0 ai ei n∈N is a Cauchy sequence in H. Hence, by completeness, there exists

ψ ∈ H such that
X∞
ψ= ai ei
i=0
and we have U ψ = {an }n∈N , so U is surjective. Moreover, we have
X ∞ ∞
X
hψ|ϕiH = hei |ψiH ei
hej |ϕiH ej
i=0 j=0 H
∞ X
X ∞
= hei |ψiH hej |ϕiH hei |ej iH
j=0 j=0
X∞ X ∞
= hei |ψiH hej |ϕiH δij
j=0 j=0
X∞
= hei |ψiH hei |ϕiH
j=0

=: {hen |ψiH }n∈N {hen |ϕiH }n∈N `2
=: hU ψ|U ϕi`2 .

Hence, U preserves the inner product, and it is therefore a unitary map.

– 40 –
4 Projectors, bras and kets

4.1 Projectors
Projectors play a key role in quantum theory, as you can see from Axioms 3 and 5.
Definition. Let H be a separable Hilbert space. Fix a unit vector e ∈ H (that is, kek = 1)
and let ψ ∈ H. The projection of ψ to e is
ψ := he|ψie
while the orthogonal complement of ψ is
ψ⊥ := ψ − ψ .
We can extend these definitions to a countable orthonormal subset {ei }i∈N ⊂ H, i.e. a
subset of H whose elements are pairwise orthogonal and have unit norm. Note that {ei }i∈N
need not be a basis of H.
Proposition 4.1. Let ψ ∈ H and let {ei }i∈N ⊂ H be an orthonormal subset. Then
(a) we can write ψ = ψ + ψ⊥ , where
∞
X
ψ := hei |ψiei , ψ⊥ := ψ − ψ
i=0
and we have
∀ i ∈ N : hψ⊥ |ei i = 0.

(b) Pythagoras’ theorem holds:

kψk2 = kψ k2 + kψ⊥ k2 .
Note that this is an extension to the finite-dimensional case.

(c) for any γ ∈ span{ei | i ∈ N}, we have the estimate

kψ − γk ≥ kψ⊥ k
with equality if, and only if, γ = ψ .
Proof. First consider the case of a finite orthonormal subset {e0 , . . . , en } ⊂ H.
(a) Let ψ and ψ⊥ be defined as in the proposition. Then ψ + ψ⊥ = ψ and
Xn

hψ⊥ |ei i = ψ − hej |ψiej ei
j=0
n
X
= hψ|ei i − hej |ψihej |ei i
j=0
n
X
= hψ|ei i − hψ|ej iδji
j=0
= hψ|ei i − hψ|ei i
=0
for all 0 ≤ i ≤ n.

– 41 –
(b) From part (a), we have
n Xn
X
hψ⊥ |ψ i = ψ⊥
hei |ψiei = hei |ψihψ⊥ |ei i = 0.
i=0 i=0

Hence, by (the finite-dimensional) Pythagoras’ theorem

kψk2 = kψ + ψ⊥ k2 = kψ k2 + kψ⊥ k2 .
Pn
(c) Let γ ∈ span{ei | 0 ≤ i ≤ n}. Then γ = i=0 γi ei for some γ0 , . . . , γn ∈ C. Hence

kψ − γk2 = kψ⊥ + ψ − γk2

and thus kψ − γk ≥ kψ⊥ k since |hei |ψi − γi |2 ≥ 0 for all 0 ≤ i ≤ n. Moreover, we

have equality if, and only if, |hei |ψi − γi | = 0 for all 0 ≤ i ≤ n, that is γ = ψ .

To extend this to a countably infinite orthonormal set {ei }i∈N , note that by part (b)
and Bessel’s inequality, we have

n n
X 2 X
|hei |ψi|2 ≤ kψk2 .

hei |ψiei =

i=0 i=0
Pn
Since |hei
|ψi|2 ≥ 0, the sequence of partial sums i=0 |hei |ψi| n∈N is monotonically
2

increasing and bounded from above by kψk. Hence, it converges and this implies that
∞
X
ψ := hei |ψiei
i=0

exists as an element of H. The extension to the countably infinite case then follows by
continuity of the inner product.

4.2 Closed linear subspaces

We will often be interested in looking at linear subspaces of a Hilbert space H, i.e. subsets
M ⊆ H such that
∀ ψ, ϕ ∈ M : ∀ z ∈ C : zψ + ϕ ∈ M.
Note that while every linear subspace M ⊂ H inherits the inner product on H to become
an inner product space, it may fail to be complete with respect to this inner product. In
other words, not every linear subspace of a Hilbert space is necessarily a sub-Hilbert space.
The following definitions are with respect to the norm topology on a normed space and
can, of course, be given more generally on an arbitrary topological space.

– 42 –
Definition. Let H be a normed space. A subset M ⊂ H is said to be open if

∀ ψ ∈ M : ∃ r > 0 : ∀ ϕ ∈ H : kψ − ϕk < r ⇒ ϕ ∈ M.

Equivalently, by defining the open ball of radius r > 0 and centre ψ ∈ H

Br (ψ) := {ϕ ∈ H | kψ − ϕk < r},

we can define M ⊂ H to be open if

∀ ψ ∈ M : ∃ r > 0 : Br (ψ) ⊆ M.

Definition. A subset M ⊂ H is said to be closed if its complement H \ M is open.

Proposition 4.2. A closed subset M of a complete normed space H is complete.

Proof. Let {ψn }n∈N be a Cauchy sequence in the closed subset M. Then, {ψn }n∈N is also
a Cauchy sequence in H, and hence it converges to some ψ ∈ H since H is complete. We
want to show that, in fact, ψ ∈ M. Suppose, for the sake of contradiction, that ψ ∈/ M,
i.e. ψ ∈ H \ M. Since M is closed, H \ M is open. Hence, there exists r > 0 such that

∀ ϕ ∈ H : kϕ − ψk < r ⇒ ϕ ∈ H \ M.

However, since ψ is the limit of {ψn }n∈N , there exists N ∈ N such that

∀ n ≥ N : kψn − ψk < r.

Hence, for all n ≥ N , we have ψn ∈ H \ M, i.e. ψn ∈/ M, contradicting the fact that

{ψn }n∈N is a sequence in M. Thus, we must have ψ ∈ M.

Corollary 4.3. A closed linear subspace M of a Hilbert space H is a sub-Hilbert space with
the inner product on H. Moreover, if H is separable, then so is M.

Knowing that a linear subspace of a Hilbert space is, in fact, a sub-Hilbert space can
be very useful. For instance, we know that there exists an orthonormal basis for the linear
subspace. Note that the converse to the corollary does not hold: a sub-Hilbert space need
not be a closed linear subspace.

4.3 Orthogonal projections

Definition. Let M ⊆ H be a (not necessarily closed) linear subspace of H. The set

M⊥ := {ψ ∈ H | ∀ ϕ ∈ M : hϕ|ψi = 0}

is called the orthogonal complement of M in H.

Proposition 4.4. Let M ⊆ H be a linear subspace of H. Then, M⊥ is a closed linear

subspace of H.

– 43 –
Proof. Let ψ1 , ψ2 ∈ M⊥ and z ∈ C. Then, for all ϕ ∈ M

hϕ|zψ1 + ψ2 i = zhϕ|ψ1 i + hϕ|ψ2 i = 0

and hence zψ1 + ψ2 ∈ M⊥ . Thus, M⊥ is a linear subspace of H. It remains to be shown

that it is also closed. Define the maps

fϕ : H → C
ψ 7→ hϕ|ψi.

Then, we can write \

M⊥ = preimfϕ ({0}).
ϕ∈M

Since the inner product is continuous (in each slot), the maps fϕ are continuous. Hence, the
pre-images of closed sets are closed. As the singleton {0} is closed in the standard topology
on C, the sets preimfϕ ({0}) are closed for all ϕ ∈ M. Thus, M⊥ is closed since arbitrary
intersections of closed sets are closed.

Remark 4.5 . Note that M⊥ is also not open (which is not necessarily the same as closed).
To clarify, equally the preimage of {0} is an open set (as {0} is open), however it is only
finite intersections of open sets that are open. So the inclusion of the intersection plays an
important role.

Note that by Pythagoras’ theorem, we have the decomposition

H = M ⊕ M⊥ := {ψ + ϕ | ψ ∈ M, ϕ ∈ M⊥ }

for any closed linear subspace M.

Proposition 4.6. For any closed linear subspace X ⊂ H its true that X ⊥⊥ = X.

Proof. Let x ∈ X, then for all y ∈ X ⊥ we have hx, yi = 0 and so x ∈ X ⊥⊥ . This gives
X ⊆ X ⊥⊥ .
Now consider z ∈ X ⊥⊥ . As X is closed from the above note we know it can be decomposed
as z = x + y for x ∈ X and y ∈ X ⊥ . We then have

0 = hy, zi
= hy, x + yi
= hy, xi + hy, yi
= kyk2
=⇒ y = 0

where the last step comes from the definiteness of k·k. So we have z ∈ X and X ⊥⊥ ⊆ X.

Proposition 4.7. For any linear subspace M ⊆ H it is true that M⊥⊥ = M, where the
latter is the topological closure of the set.

– 44 –
Proof. We start with two observations:

(i) M ⊆ M⊥⊥ , which was shown at the start of the last proof (as there was no use of
the fact that X was closed there.

(ii) If M1 ⊆ M2 then M⊥
2 ⊆ M1 , which can be shown easily.
⊥

First let’s show that M⊥⊥ ⊆ M. Clearly M ⊆ M (where the equality hold only if M is
⊥ ⊥⊥
closed), and so from (ii) we have M ⊆ M⊥ , which in turn gives M⊥⊥ ⊆ M . But M
is closed and so from the previous proposition we have M⊥⊥ ⊆ M.
Now we need to show the reverse inclusion. From Proposition 4.4 we know that M⊥⊥
is closed. Then (i) instantly tells us that M ⊆ M⊥⊥ .

Definition. Let M be a closed linear subspace of a separable Hilbert space H and fix some
orthonormal basis of M. The map

PM : H → M
ψ 7→ ψ

is called the orthogonal projector to M.

Proposition 4.8. Let PM : H → M be an orthogonal projector to M ⊆ H. Then

2 =P
(i) PM ◦ PM = PM , sometimes also written as PM M

(ii) ∀ ψ, ϕ ∈ H : hPM ψ|ϕi = hψ|PM ϕi

(iii) PM⊥ ψ = ψ⊥

(iv) PM ∈ L(H, M).

Proof. Let {ei }i∈I and {ei }i∈J be bases of M and M⊥ respectively, where I, J are disjoint
and either finite or countably infinite, such that {ei }i∈I∪J is a basis of H (Note that we
should think of I ∪ J as having a definite ordering).

(i) Let ψ ∈ H. Then

(iii) Let ψ ∈ H. Then

X X X
PM ψ + PM⊥ ψ = hei |ψiei + hei |ψiei = hei |ψiei = ψ.
i∈I i∈J i∈I∪J

Hence
PM⊥ ψ = ψ − PM ψ = ψ − ψ =: ψ⊥ .

(iv) Let ψ ∈ H. Then, by Pythagoras’ theorem,

kPM ψk kψ k kψk − kψ⊥ k
sup = sup = sup ≤1<∞
ψ∈H kψk ψ∈H kψk ψ∈H kψk

Quite interesting, and heavily used, is the converse.

Theorem 4.9. Let P ∈ L(H, H) have the properties
(i) P ◦ P = P

(ii) ∀ ψ, ϕ ∈ H : hPψ|ϕi = hψ|Pϕi.

Then, the range P(H) of P is closed and
P = PP(H) .
In other words, every projector is the orthogonal projector to some closed linear subspace.

4.4 Riesz representation theorem, bras and kets

Let H be a Hilbert space. Consider again the map
fϕ : H → C
ψ 7→ hϕ|ψi.
for ϕ ∈ H. The linearity in the second argument of the inner product implies that this map
is linear. Moreover, by the Cauchy-Schwarz inequality, we have
|fϕ (ψ)| |hϕ|ψi| kϕkkψk
sup = sup ≤ sup = kϕk < ∞.
ψ∈H kψk ψ∈H kψk ψ∈H kψk
Hence, fϕ ∈ L(H, C) =: H∗ . Therefore, to every element of ϕ of H, there is associated an
element fϕ in the dual space H∗ . In fact, the converse is also true.

– 46 –
Theorem 4.10 (Riesz representation). Every f ∈ H∗ is of the form fϕ for a unique ϕ ∈ H.

Proof. First, suppose that f = 0, i.e. f is the zero functional on H. Then, clearly, f = f0
with 0 ∈ H. Hence, suppose that f 6= 0. Since, ker f := preimf ({0}) is a closed linear
subspace, we can write
H = ker f ⊕ (ker f )⊥ .
As f 6= 0, there exists some ψ ∈ H such that ψ ∈ / ker f . Hence, ker f 6= H, and thus
(ker f ) 6= {0}. Let ξ ∈ (ker f ) \ {0} and assume, w.l.o.g., that kξk = 1. Define
⊥ ⊥

ϕ := f (ξ)ξ ∈ (ker f )⊥ .

Then, for any ψ ∈ H, we have

fϕ (ψ) − f (ψ) := hϕ|ψi − f (ψ)

Note that
f (f (ξ)ψ − f (ψ)ξ) = f (ξ)f (ψ) − f (ψ)f (ξ) = 0,
that is, f (ξ)ψ − f (ψ)ξ ∈ ker f . Since ξ ∈ (ker f )⊥ , we have

hξ|f (ξ)ψ − f (ψ)ξi = 0

and hence fϕ (ψ) = f (ψ) for all ψ ∈ H, i.e. f = fϕ . For uniqueness, suppose that

f = fϕ1 = fϕ2

for some ϕ1 , ϕ1 ∈ H. Then, for any ψ ∈ H,

0 = fϕ1 (ψ) − fϕ2 (ψ)

= hϕ1 |ψi − hϕ2 |ψi
= hϕ1 − ϕ2 |ψi

and hence, ϕ1 = ϕ2 by positive-definiteness.

Therefore, the so-called Riesz map

R : H → H∗
ϕ 7→ fϕ ≡ hϕ| · i

is a linear isomorphism, and H and H∗ be identified with one another as vector spaces.
This lead Dirac to suggest the following notation for the elements of the dual space

fϕ ≡ hϕ|.

– 47 –
Correspondingly, he wrote |ψi for the element ψ ∈ H. Since h · | · i is “a bracket”, the first
half h · | is called a bra, while the second half | · i is called a ket (nobody knows where the
missing c is). With this notation, we have

fϕ (ψ) ≡ hϕ|(|ψi) ≡ hϕ|ψi.

The notation makes evident the fact that, for any ϕ, ψ ∈ H, we can always consider the
inner product hϕ|ψi ∈ C as the result of applying fϕ ∈ H∗ to ψ ∈ H.
The advantage of this notation is that some formulæ become more intuitive and hence
are more easily memorised. For a concrete example, consider
∞
X
ψ= hei |ψiei
i=0

where {en }n∈N is a basis of H. This becomes

∞
X
|ψi = hei |ψi|ei i.
i=0

By allowing the scalar multiplication of kets also from the right, defined to yield the same
result as that on the left, we have
∞
X
= |ei ihei |ψi.
i=0

“Quite obviously”, we can bracket this as

X∞
= |ei ihei | |ψi,
i=0

where by “quite obviously”, we mean that we have a suppressed tensor product (see section
8 of the Lectures on the Geometric Anatomy of Theoretical Physics for more details on
tensors)
X∞
= |ei i ⊗ hei | |ψi.
i=0

Then, the sum in the round brackets is an element of H ⊗ H∗ . While H ⊗ H∗ is isomorphic

to End(H), its elements are maps H∗ × H → C. Hence, one needs to either make this
isomorphism explicit or, equivalently,
X ∞

= |ei i ⊗ hei | · , |ψi .
i=0

All of this to be able to write

∞
X
|ei ihei | = idH
i=0

But the original expression was already clear in the first place, without the need to add
hidden tensor products and extra rules. Of course, part of the appeal of this notation is
that one can intuitively think of something like |ei ihei | as a map H → H, by imagining
that the bra on the right acts on a ket in H, thereby producing a complex number which
becomes the coefficient of the remaining ket

|ei ihei | |ψi = |ei ihei |ψi = hei |ψi|ei i.

The major drawback of this notation, and the reason why we will not adopt it, is that
in many places (for instance, if we consider self-adjoint operators, or Hermitian operators)
this notation doesn’t produce inconsistencies only if certain conditions are satisfied. While
these conditions will indeed be satisfied most of times, it becomes extremely confusing to
formulate conditions on our objects by using a notation that only makes sense if the objects
already satisfy conditions.
Of course, as this notation is heavily used in physics and related applied sciences, it is
necessary to be able to recognise it and become fluent in it. But note that it does not make
things clearer. If anything, it makes things more complicated.

– 49 –
5 Measure theory

This and the next section will be a short recap of basic notions from measure theory and
Lebesgue integration. These are inescapable subjects if one wants to understand quantum
mechanics since

(i) the spectral theorem requires the notion of (projection-valued) measures

(ii) the most commonly emerging separable Hilbert space in quantum mechanics is the
space L2 (Rd ), whose definition needs the notion of Lebesgue integral.

5.1 General measure spaces and basic results

Definition. Let M be a non-empty set. A collection Σ ⊆ P(M ) of subsets of M is called
a σ-algebra for M if the following conditions are satisfied

(i) M ∈ Σ

(ii) if A ∈ Σ, then M \ A ∈ Σ
S∞
(iii) for any sequence {An }n∈N in Σ we have n=0 An ∈ Σ.

Remark 5.1 . If we relax the third condition so that it applies only to finite (rather than
countable) unions, we obtain the notion of an algebra, often called an algebra of sets in
order to distinguish it from the notion of algebra as a vector space equipped with a bilinear
product, with which it has nothing to do.

Remark 5.2 . Note that by condition (ii) and De Morgan’s laws, condition (iii) can be equiv-
alently stated in terms of intersections rather than unions. Recall that De Morgan’s laws
“interchange” unions with intersections and vice-versa under the complement operation.
That is, if M is a set and {Ai }i∈I is a collection of sets, then
[ \ \ [
M\ Ai = (M \ Ai ), M\ Ai = (M \ Ai ).
i∈I i∈I i∈I i∈I

A σ-algebra is closed under countably infinite unions (by definition) but also under
countably infinite intersections and finite unions and intersections.

Proposition 5.3. Let M be a set and let Σ be a σ-algebra on M . Let {An }n∈N be a
sequence in Σ. Then, for all k ∈ N, we have

(i) kn=0 An ∈ Σ
S

(ii) ∞
T Tk
n=0 An ∈ Σ and n=0 An ∈ Σ.

Proof. (i) Let the sequence {Bn }n∈N be defined as follows:

(
An if 0 ≤ n ≤ k
Bn =
∅ if n > k.

– 50 –
S∞
Then, {Bn }n∈N is a sequence in Σ, so n=0 Bn ∈ Σ. Hence, we have:
∞
[ k
[ ∞
[ k
[
Bn = Bn ∪ Bn = An
n=0 n=0 n=k+1 n=0

Sk
and thus n=0 An ∈ Σ.
S∞
(ii) As {An }n∈N is a sequence in Σ, so is {M \ An }n∈N and hence n=0 (M \ An ) ∈ Σ.
Thus, we also have
[ ∞
M\ (M \ An ) ∈ Σ
n=0

and since M \ (M \ An ) = An , by De Morgan’s laws, ∞ n=0 An ∈ Σ. That this holds

for finite intersections is shown by defining {Bn }n∈N as above.

Definition. A measurable space is a pair (M, Σ) where M is a set and Σ is a σ-algebra on

M . The elements of Σ are called measurable subsets of M .

Our goal is to assign volumes (i.e. measures) to subsets of a given set. Of course, we
would also like this assignment to satisfy some sensible conditions. However, it turns out
that one cannot sensibly assign volumes to any arbitrary collection of subsets of a given
set10 . It is necessary that the collection of subsets be a σ-algebra. In addition, just like
in topology openness and closeness are not properties of subsets but properties of subsets
with respect to a choice of topology, so does measurability of subsets only make sense with
respect to a choice of σ-algebra. In particular, a given subset could be measurable with
respect to some σ-algebra and not measurable with respect to some other σ-algebra.

Example 5.4 . The pair (M, P(M )) is a measurable space for any set M . Of course, just
like the discrete topology is not a very useful topology, the power set P(M ) is not a very
useful σ-algebra on M , unless M is countable.

Definition. The extended real line is R := R ∪ {−∞, +∞}, where the symbols −∞ and
+∞ (the latter often denoted simply by ∞) satisfy

∀ r ∈ R : −∞ ≤ r ≤ ∞

with strict inequalities if r ∈ R. The symbols ±∞ satisfy the following arithmetic rules

(i) ∀ r ∈ R : ±∞ + r = r ± ∞ = ±∞

(ii) ∀ r > 0 : r(±∞) = ±∞ r = ±∞

(iii) ∀ r < 0 : r(±∞) = ±∞ r = ∓∞

(iv) 0(±∞) = ±∞ 0 = 0.

Note that expressions such as ∞ − ∞ or −∞ + ∞ are not defined.

10
https://en.wikipedia.org/wiki/Non-measurable_set

– 51 –
Definition. Let (M, Σ) be a measurable space. A measure on (M, Σ) is a function

µ : Σ → [0, ∞],

where [0, ∞] := {r ∈ R | r ≥ 0}, such that

(i) µ(∅) = 0

(ii) for any sequence {An }n∈N in Σ with Ai ∩ Aj = ∅ whenever i 6= j, we have

∞
[ ∞
X
µ An = µ(An ).
n=0 n=0

A sequence {An }n∈N that satisfies the condition that Ai ∩ Aj = ∅ for all i 6= j is called
a pairwise disjoint sequence.

Remark 5.5 . Both sides of the equation in part (ii) of the definition of measure might take
the value ∞. There are two possible reasons why ∞ n=0 µ(An ) might be infinite. It could
P

be that µ(An ) = ∞ for some n ∈ N or, alternatively, it could be that µ(An ) < ∞ for all
n ∈ N but the sequence of partial sums { ni=0 µ(Ai )}n∈N , which is an increasing sequence
P

since µ is non-negative by definition, is not bounded above.

Definition. A measure space is a triple (M, Σ, µ) where (M, Σ) is a measurable space and
µ : Σ → [0, ∞] is a measure on M .

Example 5.6 . Let M = N and Σ = P(N). Define µ : Σ → [0, ∞] by:

(
n if A is a finite set with n elements
µ(A) = |A| :=
∞ if A is not a finite set

Then, by definition, µ(∅) = 0. Moreover, if {An }n∈N is a pairwise disjoint sequence in Σ

such that all but a finite number of An ’s are empty and each An has a finite number of
elements, we have:
[∞ X ∞
µ An = µ(An )
n=0 n=0

by counting elements. Otherwise, that is, if an infinite number of An ’s are non-empty or if

at least one An is infinite, then
∞
[ ∞
X
µ An =∞= µ(An )
n=0 n=0

and thus, the triple (N, P(N), µ) is a measure space. The measure µ on (N, P(N)) is called
counting measure and it is the usual measure on countable measurable spaces.

– 52 –
Proposition 5.7. Let (M, Σ, µ) be a measure space.

(i) If A0 , . . . , Ak ∈ Σ and Ai ∩ Aj = ∅ for all 0 ≤ i 6= j ≤ k, then:

k
[ k
X
µ An = µ(An )
n=0 n=0

(ii) If A, B ∈ Σ and A ⊆ B, then µ(A) ≤ µ(B)

(iii) If A, B ∈ Σ, A ⊆ B and µ(A) < ∞, then µ(B \ A) = µ(B) − µ(A).

Proof. (i) Let An = ∅ for all n > k. Then, {An }n∈N is a pairwise disjoint sequence in Σ
and hence, we have:
k
[ ∞
[ ∞
X k
X
µ An =µ An = µ(An ) = µ(An ).
n=0 n=0 n=0 n=0

(ii) We have B = A ∪ (B \ A) and A ∩ (B \ A) = ∅. Hence, by part (i),

µ(B) = µ(A ∪ (B \ A)) = µ(A) + µ(B \ A),

and since µ(B \ A) ≥ 0, we have µ(A) ≤ µ(B).

(iii) By decomposing B as above and rearranging, we immediately get

µ(B \ A) = µ(B) − µ(A).

Note, however, that this only makes sense if µ(A) < ∞, for otherwise we must also
have µ(B) = ∞ by part (ii), and then µ(B) − µ(A) would not be defined.

Proposition 5.8. Let (M, Σ, µ) be a measure space and let {An }n∈N be a sequence in Σ.

(i) If {An }n∈N is increasing, i.e. An ⊆ An+1 for all n ∈ N, then

∞
[
µ An = lim µ(An ).
n→∞
n=0

(ii) If µ(A0 ) < ∞ and {An }n∈N is decreasing, i.e. An+1 ⊆ An for all n ∈ N, then
\∞
µ En = lim µ(An ).
n→∞
n=0

We say that µ is (i) continuous from below and (ii) continuous from above.

Proof. (i) Define B0 := A0 and Bn := An \ An−1 . Then, {Bn }n∈N is a pairwise disjoint
sequence in Σ such that
n
[ ∞
[ ∞
[
B i = An and An = Bn .
i=0 n=0 n=0

– 53 –
Hence, we have
∞
[ ∞
[
µ An =µ Bn
n=0 n=0
∞
X
= µ(Bn )
n=0
n
X
= lim µ(Bi )
n→∞
i=0
n
[
= lim µ Bn
n→∞
i=0
= lim µ(An ).
n→∞

(ii) Define Bn := A0 \ An . Then, Bn ⊆ Bn+1 for all n ∈ N and thus, by part (i), we have
∞
[
µ Bn = lim µ(Bn )
n→∞
n=0
= lim (µ(A0 ) − µ(An ))
n→∞
= µ(A0 ) − lim µ(An ).
n→∞

Note that, by definition of Bn , we have

∞
[ ∞
\
Bn = A0 \ An .
n=0 n=0

Since µ(A0 ) < ∞, it follows from the previous proposition that

∞
\ ∞
\
µ(A0 ) − µ An = µ A0 \ An
n=0 n=0
∞
[
=µ Bn
n=0
= µ(A0 ) − lim µ(An ).
n→∞

Therefore, we have
∞
\
µ An = lim µ(An ).
n→∞
n=1

Remark 5.9 . Note that the result in the second part of this proposition need not be true
if µ(A0 ) = ∞. For example, consider (N, P(N), µ), where µ is the counting measure on
(N, P(N)). If An = {n, n + 1, n + 2, . . .}, then {An }n∈N is a decreasing sequence. Since
µ(An ) = ∞ for all n ∈ N, we have lim µ(An ) = ∞. On the other hand, ∞ n=0 An = ∅ and
T
T∞ n→∞
thus µ( n=0 An ) = 0.

– 54 –
Proposition 5.10. Let (M, Σ, µ) be a measure space. Then, µ is countably sub-additive.
That is, for any sequence {An }n∈N in Σ, we have
[∞ X ∞
µ An ≤ µ(An ).
n=0 n=0

Proof. (a) First, we show that µ(A ∪ B) ≤ µ(A) + µ(B) for any A, B ∈ Σ. Note that, for
any pair of sets A and B, the sets A \ B, B \ A and A ∩ B are pairwise disjoint and
their union is A ∪ B.

A B

A\B A∩B B\A

By writing A = (A \ B) ∪ (A ∩ B) and B = (B \ A) ∪ (A ∩ B) and using the additivity

and positivity of µ, we have

µ(A) + µ(B) = µ((A \ B) ∪ (A ∩ B)) + µ((B \ A) ∪ (A ∩ B))

= µ(A \ B) + 2µ(A ∩ B) + µ(B \ A)
= µ(A ∪ B) + µ(A ∩ B)
≥ µ(A ∪ B).

(b) We now extend this to finite unions by induction. Let {An }n∈N be a sequence in Σ
and suppose that
[n X n
µ Ai ≤ µ(Ai )
i=0 i=0
for some n ∈ N. Then, by part (a), we have
n+1
[ n
[
µ Ai = µ An+1 ∪ Ai
i=0 i=0
n
[
≤ µ(An+1 ) + µ Ai
i=0
n
X
≤ µ(An+1 ) + µ(Ai )
i=0
n+1
X
= µ(Ai ).
i=0

Hence, by induction on n with base case n = 1 and noting that the case n = 0 is
trivial (it reduces to µ(A0 ) = µ(A0 )), we have
[ n Xn
∀n ∈ N : µ Ai ≤ µ(Ai ).
i=0 i=0

– 55 –
(c) Let {An }n∈N be a sequence in Σ. Define Bn := ni=0 An . Then, {Bn }n∈N is an
S

increasing sequence in Σ. Hence, by continuity from above of µ, we have

[∞ [ ∞
µ An = µ Bn
n=0 n=0
= lim µ(Bn )
n→∞
n
[
= lim µ Ai
n→∞
i=0
n
X
≤ lim µ(Ai )
n→∞
i=0
∞
X
= µ(Ai )
i=0

which is what we wanted.

Definition. Let (M, Σ, µ) be a measure space. The measure µ is said to be finite if there
exists a sequence {An }n∈N in Σ such that ∞n=0 An = M and
S

∀ n ∈ N : µ(An ) < ∞.

Example 5.11 . The counting measure on (N, P(N)) is finite. To see this, define An := {n}.
Then, clearly ∞n=0 An = N and µ(An ) = |{n}| = 1 < ∞ for all n ∈ N.
S

5.2 Borel σ-algebras

We have already remarked the parallel between topologies and σ-algebras. A further sim-
ilarity stems from the fact that, just like for topologies, interesting σ-algebras are hardly
ever given explicitly, except in some simple cases. In general, they are defined implicitly by
some membership condition.
Proposition 5.12. Let M be a set and let {Σi : i ∈ I} be a collection of σ-algebras on M .
Define the set \
Σ := Σi = {A ∈ P(M ) | A ∈ Σi , ∀ i ∈ I}.
i∈I
Then, Σ is a σ-algebra on M .
Proof. We simply check that Σ satisfies the defining properties of a σ-algebra.
(i) We have M ∈ Σi for all i ∈ I and hence M ∈ Σ.

(ii) Let A ∈ Σ. Then, A ∈ Σi for all i ∈ I and, since each Σi is a σ-algebra, we also have
M \ A ∈ Σi for all i ∈ I. Hence, M \ A ∈ Σ.

(iii) Let {An }n∈N be a sequence in Σ. Then, {An }n∈N is a sequence in each Σi . Thus,
∞
[
∀i ∈ I : An ∈ Σi .
n=0
S∞
Hence, we also have n=0 An ∈ Σ.

– 56 –
Definition. Let M be a set and let E ⊆ P(M ) be a collection of subsets of M . The
σ-algebra generated by E, denoted σ(E), is the smallest σ-algebra on M containing all the
sets in E. That is,
A ∈ σ(E) ⇔ for all σ-algebras Σ on M : E ⊆ Σ ⇒ A ∈ Σ
or, by letting {Σi | i ∈ I} be the collection of σ-algebras on M such that E ⊆ Σ,
\
σ(E) := Σi .
i∈I

The set E is called a generating set for σ(E). Observe that the second characterisation
makes it manifest that σ(E) is indeed a σ-algebra on M by the previous proposition.
Theorem 5.13. Let (M, Σ) be a measurable space. Then, Σ = σ(E) for some E ⊆ P(M ).
This generating construction immediately allows us to link the notions of topology and
σ-algebra on a set M via the following definition.
Definition. Let (M, O) be a topological space. The Borel σ-algebra on (M, O) is σ(O).
Recall that a topology on M is a collection O ⊆ P(M ) of subsets of M which contains
∅ and M and is closed under finite intersections and arbitrary (even uncountable) unions.
The elements of the topology are called open sets. Of course, while there many choices of
σ-algebra on M , if we already have a topology O on M , then the associated Borel σ-algebra
is very convenient choice of σ-algebra since, as we will soon see, it induces a measurable
structure which is “compatible” with the already given topological structure.
This is, in fact, the usual philosophy in mathematics: we always let the stronger struc-
tures induce the weaker ones, unless otherwise specified. For instance, once we have chosen
an inner product on a space, we take the norm to be the induced norm, which induces a
metric, which in turn induces a topology on that space, from which we now know how to
obtain a canonical σ-algebra.
We remark that, while the Borel σ-algebra on a topological space is generated by the
open sets, in general, it contains much more that just the open sets.
Example 5.14 . Recall that the standard topology on R, denoted OR , is defined by
A ∈ OR ⇔ ∀ a ∈ A : ∃ ε > 0 : ∀ r ∈ R : |r − a| < ε ⇒ r ∈ A.
In fact, the elements of OR are at most countable unions of open intervals in R. Consider
now the Borel σ-algebra on (R, OR ). Let a < b. Then, for any n ∈ N, the interval (a − n1 , b)
is open. Hence, {(a − n1 , b)}n∈N is a sequence in σ(OR ). Since σ-algebras are closed under
countable intersections, we have
∞
\
(a − n1 , b) = [a, b) ∈ σ(OR ).
n=0
Hence, σ(OR ) contains, in addition to all open intervals, also all half-open intervals. It is not
difficult to show that it contains all closed intervals as well. In particular, since singletons
are closed, σ(OR ) also contains all countable subsets of R. In fact, it is non-trivial11 to
produce a subset of R which is not contained in σ(OR ).
11
https://en.wikipedia.org/wiki/Borel_set#Non-Borel_sets

– 57 –
5.3 Lebesgue measure on Rd
Definition. Let (M, Σ, µ) be a measure space. If A ∈ Σ is such that µ(A) = 0, then A is
called a null set or a set of measure zero.

The following definition is not needed for the construction of the Lebesgue measure.
However, since it is closely connected with that of null set and will be used a lot in the
future, we chose to present it here.

Definition. Let (M, Σ, µ) be a measure space and let P be some property or statement.
We say that P holds almost everywhere on M if

∃ Z ∈ Σ : µ(Z) = 0 and ∀ m ∈ M \ Z : P (m).

In other words, the property P is said to hold almost everywhere on M if it holds

everywhere on M except for a null subset of M .

Example 5.15 . Let (M, Σ, µ) be a measure space and let f, g : M → N be maps. We say
that f and g are almost everywhere equal, and we write f =a.e. g, if there exists a null set
Z ∈ Σ such that
∀ m ∈ M \ Z : f (m) = g(m).
The case f = g corresponds to Z = ∅.

Definition. A measure µ : Σ → [0, ∞] is said to be complete if every subset of every null

set is measurable, i.e.

∀ A ∈ Σ : ∀ B ∈ P(A) : µ(A) = 0 ⇒ B ∈ Σ.

Note that since for any A, B ∈ Σ, B ⊆ A implies µ(B) ≤ µ(A), it follows that every
subset of a null set, if measurable, must also be a null set.

Definition. Let (M, Σ, µ) be a measure space and let (M, +, ·) be a vector space. The
measure µ is said to be translation-invariant if

∀m ∈ M : ∀A ∈ Σ : A + m ∈ Σ and µ(A + m) = µ(A),

where A + m := {a + m | a ∈ A}.

Theorem 5.16. Let ORd be the standard topology on Rd . There exists a unique complete,
translation-invariant measure
λd : σ(ORd ) → [0, ∞]
such that for all ai , bi ∈ R with 1 ≤ i ≤ d and ai < bi , we have
d
d
Y
λ [a1 , b1 ) × · · · × [ad , bd ) = (bi − ai ).
i=1

Definition. The measure λd is called the Lebesgue measure on Rd .

– 58 –
The superscript d in λd may be suppressed if there is no risk of confusion. Note that
the Lebesgue measure on R, R2 and R3 coincides with the standard notions of length, area
and volume, with the further insight that these are only defined for the elements of the
respective Borel σ-algebras.

Proposition 5.17. The Lebesgue measure on R is finite.

Proof. Consider the sequence {[an , an + 1)}n∈N where

(
− 12 n if n is even
an =
1
2 (n + 1) if n is odd.

That is, {an }n∈N is the sequence (0, 1, −1, 2, −2, 3, −3, . . .).

−3 −2 −1 0 1 2 3

Clearly, we have ∞ n=0 [an , an + 1) = R. Since, for all n ∈ N, [an , an + 1) ∈ σ(OR ) and
S

λ [an , an + 1) = 1 < ∞, the Lebesgue measure λ is finite.

This can easily be generalised to show that λd is finite for all d ≥ 1.

5.4 Measurable maps

As we have remarked earlier, once we introduce a new structure, we should immediately
think about the associated structure-preserving maps. In the case of measurable spaces, a
measurable map is one that preserves the “measurability” structure.

Definition. Let (M, ΣM ) and (N, ΣN ) be measurable spaces. A map f : M → N is said

to be measurable if
∀ A ∈ ΣN : preimf (A) ∈ ΣM .

Note that this is exactly the definition of continuous map between topological spaces,
with “continuous” replaced by “measurable” and topologies replaced by σ-algebras.

Lemma 5.18. Let (M, ΣM ) and (N, ΣN ) be measurable spaces. A map f : M → N is

measurable if, and only if,
∀ A ∈ E : preimf (A) ∈ ΣM ,
where E ⊆ P(N ) is a generating set of ΣN .

Corollary 5.19. Let (M, OM ) and (N, ON ) be topological spaces. Any continuous map
M → N is measurable with respect to the Borel σ-algebras on M and N .

Recall that a map R → R is monotonic if it is either increasing or decreasing.

– 59 –
Corollary 5.20. Any monotonic map R → R is measurable with respect to the Borel σ-
algebra (with respect to OR ).

Proposition 5.21. Let (M, ΣM ), (N, ΣN ) and (P, ΣP ) be measurable spaces. If f : M → N

and g : N → P are both measurable, the so is their composition g ◦ f : M → P .

Proof. Let A ∈ ΣP . As g is measurable, we have preimg (A) ∈ ΣN . Then, since f is

measurable, it follows that

preimf (preimg (A)) = preimg◦f (A) ∈ ΣM .

Hence, g ◦ f is measurable.

Proposition 5.22. Let (M, ΣM ) and (N, ΣN ) be measurable spaces and let {fn }n∈N be a
sequence of measurable maps from M to N whose pointwise limit is f . Then, f is measur-
able.

Recall that {fn }n∈N converges pointwise to f : M → N if

∀ m ∈ M : lim fn (m) = f (m).

n→∞

This is in contrast with continuity, as pointwise convergence of a sequence of continuous

maps is not a sufficient condition for the continuity of the pointwise limit. In the case
of real or complex-valued maps, a sufficient condition is convergence with respect to the
supremum norm.

5.5 Push-forward of a measure

If we have a structure-preserving map f between two instances A and B of some structure,
and an object on A (which depends in some way on the structure), we can often use f to
induce a similar object on B. This is generically called the push-forward of that object
along the map f .

Proposition 5.23. Let (M, ΣM , µ) be a measure space, let (N, ΣN ) be a measurable space
and let f : M → N be a measurable map. Then, the map

f∗ µ : ΣN → [0, ∞]
A 7→ µ(preimf (A))

is a measure on (N, ΣN ) called the push-forward of µ along f .

That f∗ µ is a measure follows easily from the fact that µ is a measure and basic
properties of pre-images of maps, namely

preimf (A \ B) = preimf (A) \ preimf (B)

[ [
preimf Ai = preimf (Ai ).
i∈I i∈I

– 60 –
6 Integration of measurable functions

We will now focus on measurable functions M → R and define their integral on a subset
of M with respect to some measure on M , which is called the Lebesgue integral. Note
that, even if M ⊆ Rd , the Lebesgue integral of a function need not be with respect to the
Lebesgue measure.
The key application of this material is the definition of the Banach spaces of (classes
of) Lebesgue integrable functions Lp . The case p = 2 is especially important since L2 is, in
fact, a Hilbert space. It appears a lot in quantum mechanics where it is loosely referred to
as the space of square-integrable functions.

6.1 Characteristic and simple functions

Definition. Let M be a set and let A ∈ P(M ). The characteristic function of A, denoted
χA : M → R, is defined by (
1 if m ∈ A
χA (m) =
0 if m ∈/ A.
Example 6.1 . Below is the graph of χ(1,2] : R → R.

1 χ(1,2]

0 1 2 3

The following properties of characteristic functions are immediate.

Proposition 6.2. Let M be a set and let A, B ∈ P(M ). Then
(i) χ∅ = 0

(ii) χA∪B = χA + χB − χA∩B

(iii) χA∩B = χA χB

(iv) χM \A + χA = 1
where the addition and multiplication are pointwise and the 0 and 1 in parts (i) and (iv)
are the constant functions M → R mapping every m ∈ M to 0 ∈ R and 1 ∈ R, respectively.
Definition. Let M be a set. A function s : M → R is simple if s(M ) = {r1 , . . . , rn } for
some n ∈ N.
Equivalently, s : M → R is simple if there exist r1 , . . . , rn ∈ R and A1 , . . . , An ∈ P(M ),
for some n ∈ N, such that
X n
s= ri χAi .
i=1
So s is simple if it is a linear combination of characteristic functions.

– 61 –
Example 6.3 . Consider the simple function s : R → R given by s := χ[1,3] + 2χ[2,5] .

0 1 2 3 4 5

By observing the graph, we see that we can re-write s as

s = χ[1,2) + 3χ[2,3] + 2χ(3,5] .

Definition. A simple function is said to be in its standard form if

n
X
s= ri χAi ,
i=1

where Ai ∩ Aj = ∅ whenever i 6= j.
Any simple function can be written in standard form. It is clear that if s is in its
standard form, then Ai = preims ({ri }).
Proposition 6.4. Let (M, Σ) be a measurable space and let A, A1 , . . . , An ∈ P(M ). Then
(i) χA is measurable if, and only if, A ∈ Σ

(ii) if s = ni=1 ri χAi is a simple function in its standard form, then s is measurable if,
P

and only if, we have Ai ∈ Σ for all 1 ≤ i ≤ n.

Proof. (i) We have χA : M → R with possible values {0, 1}. Equipping R with its Borel
σ-algebra w.r.t. the standard topology, σ(OR ), we have:

(a) (⇒) We have

A = M \ preimχA ([0, 1)),
but [0, 1) ∈ σ(OR ) and so preimχA ([0, 1)) ∈ Σ. Then by property (ii) of a σ-
algebra E ∈ Σ.
(b) (⇐) Let α = [1, ∞), β = (−∞, 0) and γ = [0, 1]. Clearly α, β, γ ∈ σ(OR ) and
α ∪ β ∪ γ = R. Then

preimχA (α) = A,
preimχA (β) = ∅,
preimχA (γ) = M,

all of which are measurable sets on M and so χA is measurable.

– 62 –
(ii) First define χ
eAi := ri χAi , which satisfies
(
ri if m ∈ Ai
χ
eAi (m) =
0 if m ∈
/ Ai .

As s is in its standard form we know χ eAi ∩ χeAj = ∅ for all i 6= j. Combining these
two things, and defining αi = [ri , ∞), γi = [0, ri ], the result follow from part (i).

6.2 Integration of non-negative measurable simple functions

We begin with the definition of the integral of a non-negative, measurable, simple function.

Definition. Let (M, Σ, µ) be a measure space and let s : M → R be a nowhere negative,

measurable, simple function whose standard form is s = ni=1 ri χAi . Then, we define
P

Z n
X
s dµ := ri µ(Ai ).
M i=1

Note that the non-negativity condition is essential since µ takes values in [0, ∞], hence
we could have µ(Ai ) = ∞ for more that one Ai , and if the corresponding coefficients ri
have opposite signs, then we would have M s dµ = ∞ − ∞, which is not defined. For the
R

same reason, we are considering s : M → [0, ∞) rather than s : M → [0, ∞].

Example 6.5 . Consider the measure space (N, P(N), µ), where µ is the counting measure,
let f : N → R be non-negative and suppose that there exists N ∈ N such that f (n) = 0 for
all n > N . Then, we can write f as
N
X
f= f (n)χ{n} .
n=0

That is, f is a non-negative, measurable, simple function and therefore

Z N
X N
X
f dµ := f (n)µ({n}) = f (n).
N n=0 n=0

Hence, the “integral” of f over N with respect to the counting measure is just the sum.

The need for simple functions to be in their standard form, which was introduced to
avoid any potential ambiguity in the definition of their integral, can be relaxed using the
following lemma.

Lemma 6.6. Let (M, Σ, µ) be a measure space. Let s and t be non-negative, measurable,
simple functions M → R and let c ∈ [0, ∞). Then
Z Z Z
(cs + t) dµ = c s dµ + t dµ.
M M M

– 63 –
Pn
Proposition 6.7. Let (M, Σ, µ) be a measure space and let s = i=1 ri χAi be a non-
negative, measurable, simple function M → R not necessarily in its standard form. Then
Z n
X
s dµ = ri µ(Ai ).
M i=1

Corollary 6.8. Let (M, Σ, µ) be a measure space. Let s and t be non-negative, measurable,
simple functions M → R such that s ≤ t (that is, s(m) ≤ t(m) for all m ∈ M ). Then
Z Z
s dµ ≤ t dµ.
M M
Pn
Lemma 6.9. Let (M, Σ, µ) be a measure space and let s = i=1 ri χAi be a non-negative,
measurable, simple function M → R. Define the map

νs : Σ → [0, ∞]
Z
A 7→ sχA dµ,
M

where sχA is the pointwise product of s and χA . Then, νs is a measure on (M, Σ).

Proof. First, note that we have

Z n
Z X n
Z X n
X
sχA dµ = ri χAi χA dµ = ri χAi ∩A dµ = ri µ(Ai ∩ A).
M M i=1 M i=1 i=1

We now check that ν satisfies the defining properties of a measure.

Z Xn Xn
(i) νs (∅) := sχ∅ dµ = ri µ(Ai ∩ ∅) = ri µ(∅) = 0
M i=1 i=1

(ii) Let {Bj }j∈N be a pairwise disjoint sequence in Σ. Then

∞
[ Z
νs Bj = sχ(S∞ Bj ) dµ
j=0
j=0 M
n
X ∞
[
= ri µ (Ai ∩ Bj )
i=1 j=0
X∞ X n
= ri µ(Ai ∩ Bj )
j=0 i=1
∞ Z
X
= sχBj dµ
j=0 M
X∞
= νs (Bj ).
j=0

Thus, νs is a measure on (M, Σ).

– 64 –
6.3 Integration of non-negative measurable functions
As we are interested in measurable functions M → R, we need to define a σ-algebra on R.
We cannot use the Borel σ-algebra since we haven’t even defined a topology on R. In fact,
we can easily get a σ-algebra on R as follows.

Proposition 6.10. The set Σ := {A ∈ P(R) | A ∩ R ∈ σ(OR )} is a σ-algebra on R.

In other words, we can simply ignore the infinities in a subset of R and consider it to
be measurable if A \ {−∞, +∞} is in the Borel σ-algebra of R. We will always consider R
to be equipped with this σ-algebra.

Lemma 6.11. Let (M, Σ) be a measurable space and let f, g : M → R be measurable. Then,
the following functions are measurable.

(i) cf + g, for any c ∈ R

(ii) |f | and f 2

(iii) f g (pointwise product) if f and g are nowhere infinite

(iv) max(f, g) (defined pointwise).

Definition. Let (M, Σ, µ) be a measure space and let f : M → R be a non-negative,

measurable function. Denote by S the set of all non-negative, measurable, simple functions
s : M → R such that s ≤ f . Then, we define
Z Z
f dµ := sup s dµ.
M s∈S M

Remark 6.12 . It is often very convenient to introduce the notation

Z Z
f (x) µ(dx) ≡ f dµ,
M M

where x is a dummy variable and could be replaced by any other symbol. The reason why
this is a convenient notation is that, while some functions have standard symbols but cannot
be easily represented by an algebraic expression (e.g. characteristic functions), others are
easily expressed in terms of an algebraic formula but do not have a standard name. For
instance, it is much easier to just write
Z
x2 µ(dx)
R

than having to first denote the function R → R, x 7→ x2 by a generic f or, say, the more
specific sqR , and then write Z
sqR dµ.
R
In computer programming, this is akin to defining anonymous functions.

– 65 –
Definition. Let (M, Σ, µ) be a measure space and let f : M → R be a non-negative,
measurable function. For any A ∈ Σ (that is, any measurable subset of M ), we define
Z Z
f dµ := f χA dµ.
A M

Note that the product f χA is measurable by part (iii) of Lemma 6.11.

Lemma 6.13. Let (M, Σ, µ) be a measure space, let f, g : M → R be non-negative, mea-

surable functions such that f ≤ g, and let A, B ∈ Σ be such that A ⊆ B. Then
Z Z
(i) f dµ ≤ g dµ
M M
Z Z
(ii) f dµ ≤ f dµ.
A B

Proof. (i) Denote by Sf and Sg the sets of non-negative, measurable, simple functions
that are less than or equal to f and g, respectively. As f ≤ g, we have Sf ⊆ Sg and
hence Z Z Z Z
f dµ := sup s dµ ≤ sup s dµ =: g dµ.
M s∈Sf M s∈Sg M M

(ii) Since A ⊆ B, for any m ∈ M we have

f (m)χA (m) ≤ f (m)χB (m).

In fact, we have equality whenever m ∈ A or m ∈ M \ B, while for m ∈ B \ A the

left hand side is zero and the right-hand side is non-negative. Hence, f χA ≤ f χB and
thus, by part (i), we have
Z Z Z Z
f dµ := f χA dµ ≤ f χB dµ =: f dµ
A M M B

Proposition 6.14 (Markov inequality). Let (M, Σ, µ) be a measure space and let f : M → R
be a non-negative, measurable function. For any z ∈ [0, ∞], we have
Z
f dµ ≥ z µ(preimf ([z, ∞])).
M

Equality is achieved whenever z is an upper bound for f .

preimf ([z, ∞])

– 66 –
The following is the pivotal theorem of Lebesgue integration.

Theorem 6.15 (Monotone convergence theorem). Let (M, Σ, µ) be a measure space and let
{fn }n∈N be a sequence of non-negative, measurable functions M → R such that fn+1 ≥ fn
for all n ∈ N. If there exists a function f : M → R such that

∀ m ∈ M : lim fn (m) = f (m)

n→∞

(i.e. f is the pointwise limit of {fn }n∈N ), then f is measurable and

Z Z
lim fn dµ = f dµ.
n→∞ M M

Remark 6.16 . Observe that this result is in stark contrast with what one may be used
from Riemann integration, where pointwise converge of a sequence of integrable functions
{fn }n∈N is not a sufficient condition for the integral of the limit f to be equal to the limit
of the integrals of fn or, in fact, even for f to be integrable. For these, we need stronger
conditions on the sequence {fn }n∈N , such as uniform converge.

The definition of the integral as a supremum is clear and geometrically reasonable.

However, it is in general very difficult to evaluate the integral of any particular function
using it. The monotone convergence theorem provides a much simpler way to evaluate the
integral. One can show that, for any non-negative, measurable function f , there exists an
increasing sequence {sn }n∈N of non-negative, measurable, simple functions (which can be
explicitly constructed from f ) whose pointwise limit is f , and hence we have
Z Z
f dµ = lim sn dµ,
M n→∞ M

where the right-hand side can usually be evaluated fairly easily.

Example 6.17 . Consider the measure space (N, P(N), µ), where µ is the counting measure,
and let f : N → R be non-negative. Note that the choice of σ-algebra P(N) on N makes
every function on N (to any measurable space) measurable. Define, for every n ∈ N,
n
X
sn = f (i)χ{i} .
i=0

Then, {sn }n∈N is an increasing sequence of non-negative, measurable, simple functions

whose pointwise limit is f and therefore, by the monotone convergence theorem,
Z Z n
X ∞
X
f dµ = lim sn dµ = lim f (i)µ({i}) = f (i).
N n→∞ N n→∞
i=0 i=0

If you ever wondered why series seem to share so many properties with integrals, the reason
is that series are just integrals with respect to a discrete measure.

The monotone convergence theorem can be used to extend some of the properties of in-
tegrals of non-negative, measurable simple functions to non-negative, measurable functions
which are not-necessarily simple.

– 67 –
Lemma 6.18. Let (M, Σ, µ) be a measure space, let f, g : M → R be non-negative, mea-
surable functions and let c ∈ [0, ∞). Then
Z Z Z
(i) (cf + g) dµ = c f dµ + g dµ
M M M
Z
(ii) the map νf : Σ → [0, ∞] defined by νf (A) := f dµ is a measure on (M, Σ)
A
Z Z Z
(iii) for any A ∈ Σ, we have f dµ = f dµ + f dµ.
M A M \A

Proof. (i) Let {sn }n∈N and {tn }n∈N be increasing sequences of non-negative, measurable,
simple functions whose pointwise limits are f and g, respectively. Then, it is easy to
see that {csn + tn }n∈N is an increasing sequence of non-negative, measurable, simple
functions whose pointwise limit is cf + g. Hence, by Lemma 6.6 and the monotone
converge theorem
Z Z
(cf + g) dµ = lim (csn + tn ) dµ
M n→∞ M
Z Z
= lim c sn dµ + tn dµ
n→∞
Z M M Z

= c lim sn dµ + lim tn dµ
n→∞ M n→∞ M
Z Z
=c f dµ + g dµ.
M M

(ii) To check that νf is a measure on (M, Σ), first note that we have
Z Z
νf (∅) = f dµ := f χ∅ dµ = 0.
∅ M

Let {Ai }i∈N be a pairwise disjoint sequence in Σ. Define, for any n ∈ N,

fn := f χ(Sn Ai ) .
i=0

Since, for all n ∈ N, we have ni=0 Ai ⊆ n+1i=0 Ai and f is non-negative, {fn }n∈N is
S S

an increasing sequence of non-negative, measurable, simple functions whose pointwise

– 68 –
limits is f χ(S∞ Ai ) . Hence, by recalling Proposition 6.2, we have
i=0

∞
[ Z
νf Ai := S∞ f dµ
i=0 i=0 Ai
Z
= f χ(S∞ Ai ) dµ
i=0
MZ
= lim f χ(Sn Ai ) dµ
n→∞ M i=0
Z X n
= lim f χAi dµ
n→∞ M
i=0
Xn Z
= lim f χAi dµ
n→∞ M
i=0
∞
X
= νf (Ai ).
i=0

(iii) Note that A ∩ (M \ A) = ∅. Hence, by using the fact that νf from part (ii) is a
measure on (M, Σ), we have
Z Z Z
f dµ = νf (M ) = νf (A) + νf (M \ A) = f dµ + f dµ.
M A M \A

Part (i) of the previous lemma and the monotone convergence theorem also imply that,
for any sequence {fn }n∈N of non-negative, measurable functions, we have
∞
Z X ∞ Z
X
fn dµ = fn dµ.
M n=0 n=0 M

Again, note that this result does not hold for the Riemann integral unless stronger condi-
tions are places on the sequence {fn }n∈N .
Finally, we have a simple but crucial result for Lebesgue integration.

Theorem 6.19. Let (M, Σ, µ) be a measure space and let f : M → R be a non-negative,

measurable function. Then
Z
f dµ = 0 ⇔ f =a.e. 0.
M

Proof. (⇒) Suppose that f dµ = 0. Define An := {m ∈ M | f (m) > 1

and let
R
M n+1 }

1
sn := n+1 χAn .

By definition of An , we clearly have sn ≤ f for all n ∈ N. Hence,

Z Z
0≤ sn dµ ≤ f dµ = 0
M M

– 69 –
and thus, sn dµ = 0 for all n ∈ N. Since by definition
R
M
Z
1
sn dµ = n+1 µ(An ),
M

we must also have µ(An ) = 0 for all n ∈ N. Let A := {m ∈ M | f (m) 6= 0}. Then, as
f is non-negative, we have
∞
[ ∞
[
1
A= An = {m ∈ M | f (m) > n+1 }
n=0 n=0

and, since An ⊆ An+1 for all n ∈ N, we have

∞
[
µ(A) = µ An = lim µ(An ) = 0.
n→∞
n=0

Thus, f is zero except on the null set A. That is, f =a.e. 0.

(⇐) Suppose that f =a.e. 0. Let S be the set of non-negative, measurable, simple functions
s such that s ≤ f . As f =a.e. 0, we have s =a.e. 0 for all s ∈ S. Thus, if
n
X
s= ri χAi ,
i=1

we must have either ri = 0 or µ(Ai ) = 0 for all 1 ≤ i ≤ n. Hence, for all s ∈ S,

Z n
X
s dµ := ri µ(Ai ) = 0.
M i=1

Therefore, we have Z Z
f dµ := sup s dµ = 0.
M s∈S M

This means that, for the purposes of Lebesgue integration, null sets can be neglected
as they do not change the value of an integral. The following are some examples of this.

Corollary 6.20. Let (M, Σ, µ) be a measure space, let A ∈ Σ and let f, g : M → R be

non-negative, measurable functions.
Z
(i) If µ(A) = 0, then f dµ = 0
A
Z Z
(ii) If f =a.e. g, then f dµ = g dµ
M M
Z Z
(iii) If f ≤a.e. g, then f dµ ≤ g dµ.
M M

Proof. (i) Clearly, f χA =a.e. 0 and hence

Z Z
f dµ = f χA dµ = 0.
A M

– 70 –
(ii) As f =a.e. g, we have f − g =a.e. 0 and thus
Z Z Z
0= (f − g) dµ = f dµ − g dµ.
M M M

(iii) Let B := {m ∈ M | f (m) > g(m)}. As f ≤a.e. g, we have µ(B) = 0 and f ≤ g on

M \ B. Thus
Z Z Z Z Z
f dµ = f dµ + f dµ ≤ g dµ ≤ g dµ
M B M \B M \B M

where we used part (i) and Lemma 6.13.

Example 6.21 . Consider (R, σ(OR ), λ) and let f : R → R be the Dirichlet function
(
1 if r ∈ Q
f (r) :=
0 if r ∈ R \ Q.

The Dirichlet function is the usual example of a function which is not Riemann integrable
(on any real interval). We will now show that we can easily assign a numerical value to
its integral on any measurable subset of R. First, note that a set A ∈ σ(OR ) is null with
respect to the Lebesgue measure if, and only if,
∞
[ ∞
X
∀ ε > 0 : ∃ {In }n∈N : A⊂ In and λ(In ) < ε
n=0 n=0

where {In }n∈N is a sequence of real intervals. From this, it immediately follows that any
countable subset of R has zero Lebesgue measure. Thus, λ(Q) = 0 and hence, f =a.e. 0.
Therefore, by the previous lemmas, we have
Z
f dλ = 0
A

for any measurable subset A of R.

6.4 Lebesgue integrable functions

Since the difference ∞ − ∞ is not defined, we cannot integrate all measurable functions.
There is, however, a very small extra condition (beyond measurability) that determines the
class of functions to which we can extend our previous definition.

Definition. Let (M, Σ, µ) be a measure space and let f : M → R. The function f is said
to be (Lebesgue) integrable if it is measurable and
Z
|f | dµ < ∞.
M

We denote the set of all integrable functions M → R by LR1 (M, Σ, µ), or simply L 1 (M ) if
there is no risk of confusion.

– 71 –
For any f : M → R, we define f + := max(f, 0) and f − := max(−f, 0), which are
measurable whenever f is measurable by part (iv) of Lemma 6.11.

f f+ f−

Observe that f = f + − f − and |f | = f + + f − . Clearly, we have f + ≤ |f | and f − ≤ |f |,

and hence
Z Z Z
|f | dµ < ∞ ⇔ f dµ < ∞ and
+
f − dµ < ∞.
M M M

Definition. Let (M, Σ, µ) be a measure space and let f : M → R be integrable. Then, the
(Lebesgue) integral of f over M with respect to µ is
Z Z Z
f dµ := f + dµ − f − dµ.
M M M

It should be clear that the role of the integrability condition M |f | dµ < ∞ is to prevent
R

the integral of f from being ∞ − ∞, which is not defined.

In quantum mechanics, we usually deal with complex functions. The extension to the
complex case of the integration theory presented thus far is straightforward.

Definition. Let (M, Σ, µ) be a measure space. A complex function f : M → C is said to

be integrable if the real functions Re(f ) and Im(f ) are measurable and
Z
|f | dµ < ∞,
M

where |f | denotes the complex modulus, i.e. |f |2 = Re(f )2 + Im(f )2 . We denote the set
of all integrable complex functions by LC1 (M, Σ, µ), or simply L 1 (M ) if there is no risk of
confusion.

Definition. Let (M, Σ, µ) be a measure space and let f : M → C be integrable. We define

Z Z Z
f dµ := Re(f ) dµ + i Im(f ) dµ.
M M M

The following lemma gives the properties expected of sums and scalar multiples of
integrals. Note, however, that before we show that, say, the integral of a sum is the sum of
the integrals, it is necessary to first show that the sum of two functions in L 1 (M ) is again
in L 1 (M ).

– 72 –
Lemma 6.22. Let (M, Σ, µ) be a measure space, let f, g ∈ L 1 (M ) and let c ∈ R. Then
Z Z
(i) |f | ∈ L (M ) and f dµ ≤
1

|f | dµ
M M
Z Z
(ii) cf ∈ L 1 (M ) and cf dµ = c f dµ
M M
Z Z Z
(iii) f + g ∈ L 1 (M ) and (f + g) dµ = f dµ + g dµ
M M M

(iv) L 1 (M ) is a vector space.

Proof. (i) As |f | = |f |, we have |f | ∈ L 1 (M ). Then, by the triangle inequality,

Z Z Z
f dµ = f + dµ − −

f dµ

M ZM Z M
+ −

≤ f dµ + f dµ

ZM Z M
= +
f dµ + f − dµ
ZM M
+ −
= (f + f ) dµ
ZM
= |f | dµ.
M

(ii) We have Z Z Z
|cf | dµ = |c||f | dµ = |c| |f | dµ < ∞
M M M

and hence, we have cf ∈ L 1 (M ).

Suppose c ≥ 0. Then, (cf )+ = cf + and (cf )− = cf − and thus
Z Z Z
cf dµ = (cf ) dµ − (cf )− dµ
+
M Z M Z M
= cf + dµ − cf − dµ
M
Z M
Z
=c +
f dµ − c f − dµ
ZM M

= c (f + − f − ) dµ
ZM
=c f dµ
M

– 73 –
Now suppose c = −1. Then, (−f )+ = f − and (−f )− = f + . Thus
Z Z Z
(−f ) dµ = (−f ) dµ − (−f )− dµ
+
M ZM Z M
= f − dµ − f + dµ
MZ M Z
+ −
=− f dµ − f dµ
Z M M

= − (f + − f − ) dµ
ZM
=− f dµ.
M

The case c < 0 follows by writing c = (−1)(−c) and applying the above results.

(iii) By the triangle inequality, we have |f + g| ≤ |f | + |g| and thus

Z Z Z
|f + g| dµ ≤ |f | dµ + |g| dµ < ∞.
M M M

Hence, f + g ∈ L 1 (M ). Moreover, we have

f + g = f + − f − + g + − g − = (f + + g + ) − (f − + g − ),

where f + + g + and f − + g − are non-negative and measurable. Note that, while

(f + g)+ 6= f + + g + and (f + g)− 6= f − + g − , one can show that f + + g + and f − + g −
give an equivalent splitting of f + g to define its integral. Therefore
Z Z Z
(f + g) dµ = (f + g ) dµ − (f − + g − ) dµ
+ +
M ZM Z M Z Z
= +
f dµ + +
g dµ − f − dµ − g − dµ
ZM M Z M M
+ − + −
= (f − f ) dµ + (g − g ) dµ
ZM Z M

= f dµ + g dµ.
M M

(iv) The set of all functions from M to R is a vector space. By parts (ii) and (iii), we have
that L 1 (M ) is a vector subspace of this vector space and hence, a vector space in its
own right.

Some properties of the integrals of non-negative, measurable functions easily carry over
to general integrable functions.

Lemma 6.23. Let (M, Σ, µ) be a measure space and let f, g ∈ L 1 (M ).

Z Z
(i) If f =a.e. g, then f dµ = g dµ
M M
Z Z
(ii) If f ≤a.e. g, then f dµ ≤ g dµ.
M M

– 74 –
Just as the monotone convergence theorem was very important for integrals of non-
negative, measurable functions, there is a similar theorem that is important for integrals of
functions in L 1 (X).
Theorem 6.24 (Dominated convergence theorem). Let (M, Σ, µ) be a measure space and
let {fn }n∈N be a sequence of measurable functions which converges almost everywhere to a
measurable function f . If there exists g ∈ L 1 (M ) such that |fn | ≤a.e. g for all n ∈ N, then
(i) f ∈ L 1 (M ) and fn ∈ L 1 (M ) for all n ∈ N
Z
(ii) lim |fn − f | dµ = 0
n→∞ M
Z Z
(iii) lim fn dµ = f dµ.
n→∞ M M
Remark 6.25 . By “{fn }n∈N converges almost everywhere to f ” we mean, of course, that
there exists a null set A ∈ Σ such that
∀ m ∈ M \ A : lim fn (m) = f (m).
n→∞

6.5 The function spaces Lp (M, Σ, µ)

While in quantum mechanics we need the theory of Lp spaces only for p = 2, it is worthwhile
to develop the theory in more generality, by considering Lp for all 1 ≤ p ≤ ∞.
Definition. Let (M, Σ, µ) be a measure space and let p ∈ [1, ∞). We define
Z
p
LR (M, Σ, µ) := f : M → R f is measurable and p
|f | dµ < ∞

M
and, similarly,
Z
LCp (M, Σ, µ) := f : M → C Re(f ) and Im(f ) are measurable and p
|f | dµ < ∞ .

M
Whenever there is no risk of confusion, we lighten the notation to just L p .
Definition. Let (M, Σ, µ) be a measure space and let f : M → R. The essential supremum
of f is defined as
ess sup f := inf{c ∈ R | f ≤a.e. c}.
Then, f is said to be almost everywhere bounded (from above) if ess sup f < ∞.
Alternatively, f : M → R is almost everywhere bounded if there exists a null set A ∈ Σ
such that f restricted to M \ A is bounded.
Definition. Let (M, Σ, µ) be a measure space. We define

∞
LR (M, Σ, µ) := f : M → R f is measurable and ess sup |f | < ∞

and, similarly,

LC∞ (M, Σ, µ) := f : M → C Re(f ) and Im(f ) are measurable and ess sup |f | < ∞ .

Whenever there is no risk of confusion, we lighten the notation to just L p , for p ∈ [1, ∞].

– 75 –
All the L p spaces become vector spaces once equipped with pointwise addition and
multiplication. Let us show this is detail for LC2 .

Proposition 6.26. Let (M, Σ, µ) be a measure space. Then, LC2 is a complex vector space.

Proof. The set of all functions M → C, often denoted M C , is a vector space under pointwise
addition and multiplication. Hence, it suffices to show that LC2 is a subspace of M C .

(i) Let f ∈ LC2 and z ∈ C. As |z| ∈ R, we have:

Z Z
2 2
|zf | dµ = |z| |f |2 dµ < ∞
M M

and hence, zf ∈ LC2 .

(ii) Let f, g ∈ LC2 . Note that

|f + g|2 = (f + g)(f + g)
= (f + g)(f + g)
= f f + f g + gf + gg
= |f |2 + f g + gf + |g|2 .

Moreover, as
0 ≤ |f − g|2 = |f |2 − f g − gf + |g|2 ,
we have f g + gf ≤ |f |2 + |g|2 , and thus

|f + g|2 ≤ 2|f |2 + 2|g|2 .

Therefore Z Z Z
|f + g|2 dµ ≤ 2 |f |2 dµ + 2 |g|2 dµ < ∞,
M M M

and hence f + g ∈ LC2 .

Ideally, we would like to turn all these L p space into Banach spaces. Let us begin by
equipping them with a weaker piece of extra structure.

Proposition 6.27. Let (M, Σ, µ) be a measure space and let p ∈ [0, ∞]. Then, the maps
k · kp : L p → R defined by
Z 1
p
p

|f | dµ for 1 ≤ p < ∞

kf kp := M

ess sup |f | for p = ∞


are semi-norms on L p . That is, for all z ∈ C and f, g ∈ L p ,

(i) kf kp ≥ 0

(ii) kzf kp = |z|kf kp

– 76 –
(iii) kf + gkp ≤ kf kp + kgkp .
In other words, the notion of semi-norm is a generalisation of that of norm obtained by
relaxing the definiteness condition. If the measure space (M, Σ, µ) is such that the empty
set is the only null set, then k · kp is automatically definite and hence, a norm.
Example 6.28 . Consider (N, P(N), µ), where µ is the counting measure. Then, as µ(A) is
the cardinality of A, the only null set is the empty set. Thus, recalling that functions on N
are just sequences, the maps
 1
∞
 X p
|an |p
for 1 ≤ p < ∞


k{an }n∈N kp = n=0

sup{|a | | n ∈ N} for p = ∞

n

are norms on L p (N). In particular, note that we have L 2 (N) = `2 (N).

However, in general measure spaces, we only have

kf kp = 0 ⇔ f =a.e. 0,

as we have shown in Theorem 6.19 for L 1 , and it is often very easy to produce an f 6= 0
such that kf kp = 0. The solution to this problem is to construct new spaces from the L p in
which functions that are almost everywhere equal are, in fact, the same function. In other
words, we need to consider the quotient space of L p by the equivalence relation “being
almost everywhere equal”.
Definition. Let M be a set. An equivalence relation on M is a set ∼ ⊆ M × M such that,
writing a ∼ b for (a, b) ∈ ∼, we have
(i) a ∼ a (reflexivity)

(ii) a ∼ b ⇔ b ∼ a (symmetry)

(iii) (a ∼ b and b ∼ c) ⇒ a ∼ c (transitivity)

for all a, b, c ∈ M . If ∼ is an equivalence relation on M , we define the equivalence class of
m ∈ M by [m] := {a ∈ M | m ∼ a} and the quotient set of M by ∼ by

M/∼ := {[m] | m ∈ M }.

It is easy to show that M/∼ is a partition of M , i.e.

[
M= [m] and [a] ∩ [b] = ∅ whenever a 6= b.
m∈M

In fact, the notions of equivalence relation on M and partition of M are one and the same.
Lemma 6.29. Let (M, Σ, µ) be a measure space and let ∼ be defined by

f ∼g :⇔ f =a.e. g.

Then, ∼ is an equivalence relation on L p .

– 77 –
Proof. Let f, g, h ∈ L p . Clearly, f ∼ f and f ∼ g ⇔ g ∼ f . Now suppose that f ∼ g and
g ∼ h. Then, there exist null sets A, B ∈ Σ such that f = g on M \ A and g = h on M \ B.
Recall that σ-algebras are closed under intersections and hence, A ∩ B ∈ Σ. Obviously, we
have f = h on M \ (A ∩ B) and, since

µ(A ∩ B) ≤ µ(A ∪ B) ≤ µ(A) + µ(B) = 0,

the set A ∩ B is null. Thus, f ∼ h.

Definition. Let (M, Σ, µ) be a measure space and let f ∼ g ⇔ f =a.e. g. We define

Lp := L p /∼ = {[f ] | f ∈ L p }.

Lemma 6.30. Let (M, Σ, µ) be a measure space. Then, the maps

(i) + : Lp × Lp → Lp defined by [f ] + [g] := [f + g]

(ii) · : C × Lp → Lp defined by z[f ] := [zf ]

(iii) k · kp : Lp → R defined by k[f ]kp := kf kp

are well-defined. Moreover, k · kp : Lp → R is a norm on Lp .

Lemma 6.31 (Hölder’s inequality). Let (M, Σ, µ) be a measure space and let p, q ∈ [1, ∞]
be such that p1 + 1q = 1 (where ∞
1
:= 0). Then, for all measurable functions f, g : M → C,
we have Z Z 1 Z 1
p q
p q

f g dµ ≤ |f | dµ |g| dµ .

M M M

Hence, if [f ] ∈ Lp and [g] ∈ Lq , with 1

p + 1
q = 1 , then [f g] ∈ L1 and

k[f g]k1 ≤ k[f ]kp k[g]kq .

The equality holds if and only if |f |p and |g|q are linearly dependent on L1 . That is,
there exists non-negative real numbers α, β ∈ R such that

α|f |p =a.e. β|g|q

Note it is clear that α and β cannot both vanish.

Theorem 6.32. The spaces Lp are Banach spaces for all p ∈ [0, ∞].

We have already remarked that the case p = 2 is special in that L2 is the only Lp space
which can be made into a Hilbert space.

Proposition 6.33. Let (M, Σ, µ) be a measure space and define

h·|·iL2 : L2 × L2 → C
Z
([f ], [g]) 7→ f g dµ.
M

Then, h·|·iL2 is well-defined and it is a sesqui-linear inner product on L2 .

– 78 –
Proof. First note that if [f ] ∈ L2 , then [f ] ∈ L2 and hence, by Hölder inequality, [f g] ∈ L1 .
This ensures that [f ][g] L2 ∈ C for all [f ], [g] ∈ L2 .

To show well-definedeness, let f 0 =a.e. f and g 0 =a.e. g. Then, f 0 g 0 =a.e. f g and thus
Z Z

0 0 0

[f ] [g ] L2 :=
0
f g dµ = f g dµ =: [f ][g] L2 .
M M

Now, let [f ], [g], [h] ∈ L2 and z ∈ C. Then

Z Z
(i) [f ] [g] L2 =

f g dµ = f g dµ = [g][f ] L2
M M

(ii) We have
Z

[f ]z[g] + [h] L2 = f (zg + h) dµ
M
Z Z
=z f g dµ + f h dµ

M
M
= z [f ][g] L2 + [f ][h] L2
Z Z
(iii) We have [f ][f ] L2 = |f |2 dµ ≥ 0 and

f f dµ =
M M
Z
|f |2 dµ = 0

[f ][f ] L2 = 0 ⇔ ⇔ k[f ]k2 = 0.
M

Thus, [f ] = 0 := [0].

The last part of the proof also shows that h·|·iL2 induces the norm k · k2 , with respect
to which L2 is a Banach space. Hence, (L2 , h·|·iL2 ) is a Hilbert space.

Remark 6.34 . The inner product h·|·iL2 on L2C (N, P(N), µ) coincides with the inner product
h·|·i`2 on `2 (N) defined in the section on separable Hilbert spaces.

– 79 –
7 Self-adjoint and essentially self-adjoint operators
While we have already given some of the following definitions in the introductory section
on the axioms of quantum mechanics, we reproduce them here for completeness.

7.1 Adjoint operators

Definition. A linear map or operator A : DA → H is said to be densely defined if DA is
dense in H, i.e.
∀ ε > 0 : ∀ ψ ∈ H : ∃ α ∈ DA : kα − ψk < ε.

Equivalently, DA = H, i.e. for every ψ ∈ H there exists a sequence {αn }n∈N in DA

whose limit is ψ.

Definition. Let A : DA → H be a densely defined operator on H. The adjoint of A is the

operator A∗ : DA∗ → H defined by

(i) DA∗ := {ψ ∈ H | ∃ η ∈ H : ∀ α ∈ DA : hψ|Aαi = hη|αi}

(ii) A∗ ψ := η.

Proposition 7.1. The adjoint operator A∗ : DA∗ → H is well-defined.

Proof. Let ψ ∈ H and let η, ηe ∈ H be such that

∀ α ∈ DA : hψ|Aαi = hη|αi and hψ|Aαi = he

η |αi.

Then, for all α in DA , we have

and hence, by positive-definiteness, η = ηe.

If A and B are densely defined and DA = DB , then the pointwise sum A + B is clearly
densely defined and hence, (A + B)∗ exists. However, we do not have (A + B)∗ = A∗ + B ∗
in general, unless one of A and B is bounded, but we do have the following result.

Proposition 7.2. If A is densely defined, then

(A + z idDA )∗ = A∗ + z idDA

for any z ∈ C.

The identity operator idDA is usually suppressed in the notation, so that the above
equation reads (A + z)∗ = A∗ + z.

Definition. Let A : DA → H be a linear operator. The kernel and range of A are

ker(A) := {α ∈ DA | Aα = 0}, ran(A) := {Aα | α ∈ DA }.

The range is also called image and im(A) and A(DA ) are alternative notations.

– 80 –
Proposition 7.3. An operator A : DA → H is

(i) injective if, and only if, ker(A) = {0}

(ii) surjective if, and only if, ran(A) = H.

Definition. An operator A : DA → H is invertible if there exists an operator B : H → DA

such that A ◦ B = idH and B ◦ A = idDA .

Proposition 7.4. An operator A is invertible if, and only if,

ker(A) = {0} and ran(A) = H.

Proposition 7.5. Let A be densely defined. Then, ker(A∗ ) = ran(A)⊥ .

Proof. We have

ψ ∈ ker(A∗ ) ⇔ A∗ ψ = 0 ⇔ ∀ α ∈ DA : hψ|Aαi = 0 ⇔ ψ ∈ ran(A)⊥ .

Definition. Let A : DA → H and B : DB → H be operators. We say that B is an extension

of A, and we write A ⊆ B, if

(i) DA ⊆ DB

(ii) ∀ α ∈ DA : Aα = Bα.

Proposition 7.6. Let A, B be densely defined. If A ⊆ B, then B ∗ ⊆ A∗ .

Proof. (i) Let ψ ∈ DB ∗ . Then, there exists η ∈ H such that

∀ β ∈ DB : hψ|Bβi = hη|βi.

In particular, as A ⊆ B, we have DA ⊆ DB and thus

∀ α ∈ DA ⊆ DB : hψ|Bαi = hψ|Aαi = hη|αi.

Therefore, ψ ∈ DA∗ and hence, DB ∗ ⊆ DA∗ .

(ii) From the above, we also have B ∗ ψ := η =: A∗ ψ for all ψ ∈ DB ∗ .

7.2 The adjoint of a symmetric operator

Definition. A densely defined operator A : DA → H is called symmetric if

∀ α, β ∈ DA : hα|Aβi = hAα|βi.

Remark 7.7 . In the physics literature, symmetric operators are usually referred to as Her-
mitian operators. However, this notion is then confused with the that of self-adjointness
when physicists say that observables in quantum mechanics correspond to Hermitian op-
erators, which is not the case. On the other hand, if one decides to use Hermitian as a
synonym of self-adjoint, it is then not true that all symmetric operators are Hermitian. In
order to prevent confusion, we will avoid the term Hermitian altogether.

– 81 –
Proposition 7.8. If A is symmetric, then A ⊆ A∗ .

Proof. Let ψ ∈ DA and let η := Aψ. Then, by symmetry, we have

∀ α ∈ DA : hψ|Aαi = hη|αi

and hence ψ ∈ DA∗ . Therefore, DA ⊆ DA∗ and A∗ ψ := η = Aψ.

Definition. A densely defined operator A : DA → H is self-adjoint if A = A∗ . That is,

(i) DA = DA∗

(ii) ∀ α ∈ DA : Aα = A∗ α.

Remark 7.9 . Observe that any self-adjoint operator is also symmetric, but a symmetric
operator need not be self-adjoint.

Corollary 7.10. A self-adjoint operator is maximal with respect to self-adjoint extension.

Proof. Let A, B be self-adjoint and suppose that A ⊆ B. Then

A ⊆ B = B ∗ ⊆ A∗ = A

and hence, B = A.

In fact, self-adjoint operators are maximal even with respect to symmetric extension,
for we would have B ⊆ B ∗ instead of B = B ∗ in the above equation.

7.3 Closability, closure, closedness

Definition. (i) A densely defined operator A is called closable if its adjoint A∗ is also
densely defined.

(ii) The closure of a closable operator is A := A∗∗ = (A∗ )∗ .

(iii) An operator is called closed if A = A.

Remark 7.11 . Note that we have used the overline notation in several contexts with different
meanings. When applied to complex numbers, it denotes complex conjugation. When
applied to subsets of a topological space, it denotes their topological closure. Finally, when
applied to (closable) operators, it denotes their closure as defined above.

Proposition 7.12. A symmetric operator is necessarily closable.

Proof. Let A be symmetric. Then, A ⊆ A∗ and hence, DA ⊆ DA∗ . Since a symmetric

operators are densely defined, we have

H = DA ⊆ DA∗ ⊆ H.

Hence, DA∗ = H.

– 82 –
Note carefully that the adjoint of a symmetric operator need not be symmetric. In
particular, we cannot conclude that A∗ ⊆ A∗∗ . In fact, the reversed inclusion holds.
Proposition 7.13. If A is symmetric, then A∗∗ ⊆ A∗ .
Proof. Since A is symmetric, we have A ⊆ A∗ . Hence, A∗∗ ⊆ A∗ by Proposition 7.6.

Lemma 7.14. For any closable operator A, we have A ⊆ A∗∗ .

Proof. Recall that DA∗ := {ψ ∈ H | ∀ α ∈ DA : hψ|Aαi = hη|αi}. Then
∀ ψ ∈ DA∗ : ∀ α ∈ DA : hψ|Aαi = hA∗ ψ|αi.
Since ψ and α above are “dummy variables”, and the order of quantifiers of the same type
is immaterial, we have
∀ ψ ∈ DA : ∀ α ∈ DA∗ : hα|Aψi = hA∗ α|ψi.
Then, taking the conjugate on both inner product terms and moving them across the equal
sign, we have
∀ ψ ∈ DA : ∀ α ∈ DA∗ : hψ|A∗ αi = hAψ|αi.
Now, letting η := Aψ ∈ H we see that ψ ∈ DA∗∗ , and so DA ⊆ DA∗∗ . Moreover, by
definition, A∗∗ ψ := η := Aψ for all ψ ∈ A, and thus A ⊆ A∗∗ .

Corollary 7.15. If A is symmetric, then A ⊆ A ⊆ A∗ .

Corollary 7.16. If A is symmetric, then A is symmetric.
Theorem 7.17. Let A be a densely defined operator.
(i) The operator A∗ is closed

(ii) If A is invertible, we have (A−1 )∗ = (A∗ )−1

−1
(iii) If A is invertible and closable and A is injective, then A−1 = A .
These theorems are proved using a graphical formulation which has not been presented
here, so in order to save space all the proofs are not given. However as a example of the
technique, we shall provide the proof for (i).12

Proof. Let H and K be two Hilbert spaces and let A : H → K be densely defined. We
define the graph of A as
Γ(A) := {(h, Ah) | h ∈ DH }.
In this context, we say A is closed if and only if Γ(A) is closed w.r.t. the product topology
on H ⊕ K. Next define the operator
J : H⊕K →H⊕K
h ⊕ k 7→ (−k) ⊕ h.
We now wish to show that Γ(A∗ ) ∼ = [J (Γ(A))]⊥ , as the right hand side is closed due
to the orthogonal projection, which gives us that A∗ is closed.
12
Many thanks to Alfredo Sepulveda-Ximenez for providing a brilliant answer to this on Quora.

– 83 –
(⇒) Let x ∈ DA and y ∈ DA∗ . We have y ⊕ A∗ (y) ∈ Γ(A∗ ), and

hy ⊕ A∗ (y), J (x ⊕ A(x)i = hy ⊕ A∗ (y), −A(x) ⊕ xi

= −hy, A(x)i + hA∗ (y), xi
= −hy, A(x)i + hy, A(x)i
=0

and so Γ(A∗ ) ⊆ [J (Γ(A))]⊥ .

(⇐) Let x ∈ DA and y ⊕ z ∈ [J (Γ(A))]⊥ . Then,

hy ⊕ z, −A(x) ⊕ xi = 0
∴ hy, A(x)i = hz, xi,

which, from the definition of the adjoint, tells us that y ∈ DA∗ and z = A∗ (y) and so
y ⊕ z ∈ Γ(A∗ ) and [J (Γ(A))]⊥ ⊆ Γ(A∗ ).

7.4 Essentially self-adjoint operators

Usually, checking that an operator is symmetric is easy. By contrast, checking that an
operator is self-adjoint (directly from the definition) requires the construction of the adjoint,
which is not always easy. However, since every self-adjoint operator is symmetric, we
can first check the symmetry property, and then determine criteria for when a symmetric
operator is self-adjoint, or for when a self-adjoint extension exists.
Two complications with the extension approach are that, given a symmetric operator,
there could be no self-adjoint extension at all, or there may be several different self-adjoint
extensions. There is, however, a class of operators for which the situation is much nicer.

Definition. A symmetric operator A is called essentially self-adjoint if A is self-adjoint.

This is weaker than the self-adjointness condition.

Proposition 7.18. If A is self-adjoint, then it is essentially self-adjoint.

Proof. If A = A∗ , then A ⊆ A∗ and A∗ ⊆ A. Hence, A∗∗ ⊆ A∗ and A∗ ⊆ A∗∗ , so A∗ = A∗∗ .

∗
Similarly, we have A∗∗ = A∗∗∗ , which is just A = A .

Theorem 7.19. If A is essentially self-adjoint, then there exists a unique self-adjoint ex-
tension of A, namely A.

Proof. (i) Since A is symmetric, it is closable and hence A exists.

(ii) By Lemma 7.14, we have A ⊆ A, so A is an extension of A.

(iii) Suppose that B is another self-adjoint extension of A. Then, A ⊆ B = B ∗ , hence

B ∗∗ ⊆ A∗ , and thus A∗∗ ⊆ B ∗∗∗ = B. This means that A ⊆ B, i.e. B is a self-adjoint
extension of the self-adjoint operator A. Hence, B = A by Corollary 7.10.

– 84 –
Remark 7.20 . One may get the feeling at this point that checking for essential self-adjointness
of an operator A, i.e. checking that A∗∗ = A∗∗∗ , is hardly easier than checking whether A
is self-adjoint, that is, whether A = A∗ . However, this is not so. While we will show below
that there is a sufficient criterion for self-adjointness which does not require to calculate
the adjoint, we will see that there is, in fact, a necessary and sufficient criterion to check
for essential self-adjointness of an operator without calculating a single adjoint.
Remark 7.21 . If a symmetric operator A fails to even be essentially self-adjoint, then there
is either no self-adjoint extension of A or there are several.
Definition. Let A be a densely defined operator. The defect indices of A are

d+ := dim(ker(A∗ − i)), d− := dim(ker(A∗ + i)),

where by A∗ ± i we mean, of course, A∗ ± i · idDA∗ .

Theorem 7.22. A symmetric operator has a self-adjoint extension if its defect indices
coincide. Otherwise, there exist no self-adjoint extension.
Remark 7.23 . We will later see that if d+ = d− = 0, then A is essentially self-adjoint.

7.5 Criteria for self-adjointness and essential self-adjointness

Theorem 7.24. A symmetric operator A is self-adjoint if (but not only if )

∃ z ∈ C : ran(A + z) = H = ran(A + z).

Proof. Since A is symmetric, by Proposition 7.8, we have A ⊆ A∗ . Hence, it remains to be

shown that A∗ ⊆ A. To that end, let ψ ∈ DA∗ and let z ∈ C. Clearly,

A∗ ψ + zψ ∈ H.

Now suppose that z satisfies the hypothesis of the theorem. Then, as ran(A + z) = H,

∃ α ∈ DA : A∗ ψ + zψ = (A + z)α.

By using the symmetry of A, we have that, for any β ∈ DA ,

hψ|(A + z)βi = h(A + z)∗ ψ|βi

Since β ∈ DA was arbitrary and ran(A + z) = H, we have

∀ϕ ∈ H : hψ|ϕi = hα|ϕi.

Hence, by positive-definiteness of the inner product, we have ψ = α, thus ψ ∈ DA and

therefore, A∗ ⊆ A.

– 85 –
Theorem 7.25. A symmetric operator A is essentially self-adjoint if, and only if,

∃ z ∈ C \ R : ran(A + z) = H = ran(A + z).

The following criterion for essential self-adjointness, which does require the calculation
of A∗ ,is equivalent to the previous result and, in some situations, it can be easier to check.

Theorem 7.26. A symmetric operator A is essentially self-adjoint if, and only if,

∃ z ∈ C \ R : ker(A∗ + z) = {0} = ker(A∗ + z).

Proof. We show that this is equivalent to the previous condition. Recall that if M is a
linear subspace of H, then M⊥ is closed and hence, M⊥⊥ = M (Proposition 4.7). Thus,
by Proposition 7.5, we have

ran(A + z) = ran(A + z)⊥⊥ = ker(A∗ + z)⊥

and, similarly,
ran(A + z) = ker(A∗ + z)⊥ .
Since H⊥ = {0}, the above condition is equivalent to

ran(A + z) = H = ran(A + z).

– 86 –
8 Spectra and perturbation theory
We will now focus on the spectra of operators and on the decomposition of the spectra of
self-adjoint operators. The significance of spectra is that the axioms of quantum mechanics
prescribe that the possible measurement values of an observable (which is, in particular, a
self-adjoint operator) are those in the so-called spectrum of the operator.
A common task in almost any quantum mechanical problem that you might wish to
solve is to determine the spectrum of some observable. This is usually the Hamiltonian,
or energy operator, since the time evolution of a quantum system is governed by the expo-
nential of the Hamiltonian, which is more practically determined by first determining its
spectrum.
More often than not, it is not possible to determine the spectrum of an operator exactly
(i.e. analytically). One then resorts to perturbation theory which consists in expressing the
operator whose spectrum we want to determine as the sum of an operator whose spectrum
can be determined analytically and another whose contribution is “small” in some sense to
be made precise.

8.1 Resolvent map and spectrum

Definition. The resolvent map of an operator A is the map
RA : ρ(A) → L(H)
z 7→ (A − z)−1 ,
where L(H) ≡ L(H, H) and ρ(A) is the resolvent set of A, defined as
ρ(A) := {z ∈ C | (A − z)−1 ∈ L(H)}.
Remark 8.1 . Checking whether a complex number z belongs to ρ(A) may seem like a daunt-
ing task and, in general, it is. However, we will almost exclusively be interested in closed
operators, and the closed graph theorem states that if A is closed, then (A − z)−1 ∈ L(H)
if, and only if, A − z is bijective.
Definition. The spectrum of an operator A is σ(A) := C \ ρ(A).
Definition. A complex number λ ∈ C is said to be an eigenvalue of A : DA → H if
∃ ψ ∈ DA \ {0} : Aψ = λψ.
Such an element ψ is called an eigenvector of A associated to the eigenvalue λ.
Corollary 8.2. Let λ ∈ C be an eigenvalue of A. Then, λ ∈ σ(A).
Proof. If λ is an eigenvalue of A, then there exists ψ ∈ DA \ {0} such that Aψ = λψ, i.e.
(A − λ)ψ = 0.
Thus, ψ ∈ ker(A − λ) and hence, since ψ 6= 0, we have
ker(A − λ) 6= {0}.
This means that A − λ is not injective, hence not invertible and thus, λ ∈
/ ρ(A). Then, by
definition, λ ∈ σ(A).

– 87 –
Remark 8.3 . If H is finite-dimensional, then the converse of the above corollary holds ad
hence, the spectrum coincides with the set of eigenvalues. However, in infinite-dimensional
spaces, the spectrum of an operator contains more than just the eigenvalues of the operator.

8.2 The spectrum of a self-adjoint operator

Recall that a self-adjoint operator is necessarily closed since A = A∗ implies A = A∗∗ .
While the following refinement of the notion of spectrum can be made in greater generality,
we will primarily be interested in the case of self-adjoint operators.

Definition. Let A be a self-adjoint operator. Then, we define

(i) the pure point spectrum of A

σpp (A) := {z ∈ C | ran(A − z) = ran(A − z) 6= H}

(ii) the point embedded in continuum spectrum of A

σpec (A) := {z ∈ C | ran(A − z) 6= ran(A − z) 6= H}

(iii) the purely continuous spectrum of A

σpc (A) := {z ∈ C | ran(A − z) 6= ran(A − z) = H}.

These form a partition of σ(A), i.e. they are pairwise disjoint and their union is σ(A).

Definition. Let A be a self-adjoint operator. Then, we further define

(i) the point spectrum of A

σp (A) := σpp (A) ∪ σpec (A) = {z ∈ C | ran(A − z) 6= H}

(ii) the continuous spectrum of A

σc (A) := σpec (A) ∪ σpc (A) = {z ∈ C | ran(A − z) 6= ran(A − z)}.

Clearly, σp (A) ∪ σc (A) = σ(A) but, since σp (A) ∩ σc (A) = σpec (A) is not necessarily
empty, the point and continuous spectra do not form a partition of the spectrum in general.

Lemma 8.4. Let A be self-adjoint and let λ be an eigenvalue of A. Then, λ ∈ R.

Proof. Let ψ ∈ DA \ {0} be an eigenvector of A associated to λ. By self-adjointness of A,

λhψ|ψi = hψ|λψi = hψ|Aψi = hAψ|ψi = hλψ|ψi = λhψ|ψi.

Thus, we have
(λ − λ)hψ|ψi = 0
and since ψ 6= 0, it follows that λ = λ. That is, λ ∈ R.

– 88 –
Theorem 8.5. If A is a self-adjoint operator, then the elements of σp (A) are precisely the
eigenvalues of A.

Proof. (⇐) Suppose that λ is an eigenvalue of A. Then, by self-adjointness of A,

6 ker(A − λ) = ker(A∗ − λ) = ker((A − λ)∗ ) = ran(A − λ)⊥ = ran(A − λ)⊥ ,

{0} =

where we made use of our previous lemma. Hence, we have

ran(A − λ) = ran(A − λ)⊥⊥ 6= {0}⊥ = H

and thus, λ ∈ σp (A).

(⇒) We now need to show that if λ ∈ σp (A), then λ is an eigenvalue of A. By contraposi-

tion, suppose that λ ∈ C is not an eigenvalue of A. Note that if λ is real, then λ = λ
while if λ is not real, then λ is not real. Hence, if λ is not an eigenvalue of A, then
neither is λ. Therefore, there exists no non-zero ψ in DA such that Aψ = λψ. Thus,
we have

{0} = ker(A − λ) = ker(A∗ − λ) = ker((A − λ)∗ ) = ran(A − λ)⊥

and hence
ran(A − λ) = ran(A − λ)⊥⊥ = {0}⊥ = H.
Therefore, λ ∈
/ σp (A).

Remark 8.6 . The contrapositive of the statement P ⇒ Q is the statement ¬Q ⇒ ¬P , where

the symbol ¬ denotes logical negation. A statement and its contrapositive are logically
equivalent and “proof by contraposition” simply means “proof of the contrapositive”.

8.3 Perturbation theory for point spectra of self-adjoint operators

Before we move on to perturbation theory, we will need some preliminary definitions. First,
note that if ψ and ϕ are both eigenvectors of an operator A associated to some eigenvalue
λ, then, for any z ∈ C, the vector zψ + ϕ is either zero or it is again an eigenvector of A
associated to λ.

Definition. Let A be an operator and let λ be an eigenvalue of A.

(i) The eigenspace of A associated to λ is

EigA (λ) := {ψ ∈ DA | Aψ = λψ}.

(ii) The eigenvalue λ is said to be non-degenerate if dim EigA (λ) = 1, and degenerate if
dim EigA (λ) > 1.

(iii) The degeneracy of λ is dim EigA (λ).

Remark 8.7 . Of course, it is possible that dim EigA (λ) = ∞ in general. However, in this
section, we will only consider operators whose eigenspaces are finite-dimensional.

– 89 –
Lemma 8.8. Eigenvectors associated to distinct eigenvalues of a self-adjoint operator are
orthogonal.

Proof. Let λ, λ0 be distinct eigenvalues of a self-adjoint operator A and let ψ, ϕ ∈ DA \ {0}

be eigenvectors associated to λ and λ0 , respectively. As A is self-adjoint, we already know
that λ, λ0 ∈ R. Then, note that

(λ − λ0 )hψ|ϕi = λhψ|ϕi − λ0 hψ|ϕi

Since λ − λ0 6= 0, we must have hψ|ϕi = 0.

A. Unperturbed spectrum
Let H0 be a self-adjoint operator whose eigenvalues and eigenvectors are known and satisfy

H0 enδ = hn enδ ,

where

• the index n varies either over N or some finite range 1, 2, . . . , N

• the real numbers hn are the eigenvalues of H0

• the index δ varies over the range 1, 2, . . . , d(n), with d(n) := dim EigH0 (hn )

• for each fixed n, the set

{enδ ∈ DH0 | 1 ≤ δ ≤ d(n)}
is a linearly independent subset (in fact, a Hamel basis) of EigH0 (hn ).

Note that, since we are assuming that all eigenspaces of H0 are finite-dimensional, EigH0 (hn )
is a sub-Hilbert space of H and hence, for each fixed n, we can choose the enδ so that

henα |enβ i = δαβ .

In fact, thanks to our previous lemma, we can choose the eigenvectors of H0 so that

henα |emβ i = δnm δαβ .

Let W : DH0 → H be a not necessarily self-adjoint operator. Let λ ∈ (−ε, ε) ⊆ R and

consider the real one-parameter family of operators {Hλ | λ ∈ (−ε, ε)}, where

Hλ := H0 + λW.

Further assume that Hλ is self-adjoint for all λ ∈ (−ε, ε). Recall, however, that this
assumption does not force W to be self-adjoint.

– 90 –
We seek to understand the eigenvalue equation for Hλ ,

Hλ enδ (λ) = hnδ (λ)enδ (λ),

by exploiting the fact that it coincides with the eigenvalue equation for H0 when λ = 0.
In particular, we will be interested in the lifting of the degeneracy of hn (for some fixed
n) once the perturbation W is “switched on”, i.e. when λ 6= 0. Indeed, it is possible, for
instance, that while the two eigenvectors en1 and en2 are associated to the same (degenerate)
eigenvalue hn of H0 , the “perturbed” eigenvectors en1 (λ) and en2 (λ) may be associated to
different eigenvalues of Hλ . Hence the reason why we added a δ-index to the eigenvalue in
the above equation. Of course, when λ = 0, we have hnδ (λ) = hn for all δ.

B. Formal power series ansatz13

In order to determine hnδ (λ) and enδ (λ), we make, for both, the following ansatz
(1) (2)
hnδ (λ) =: hn + λθnδ + λ2 θnδ + O(λ3 )
(1) (2)
enδ (λ) =: enδ + λnδ + λ2 nδ + O(λ3 ),
(1) (2) (1) (2)
where θnδ , θnδ ∈ R and nδ , nδ ∈ DH0 .

Remark 8.9 . Recall that the Big O notation is defined as follows. If f and g are functions
I ⊆ R → R and a ∈ I, then we write

f (x) = O(g(x)) as x → a

to mean
∃ k, M > 0 : ∀ x ∈ I : 0 < |x − a| < k ⇒ |f (x)| < M |g(x)|.
The qualifier “as x → a” can be omitted when the value of a is clear from the context. In
our expressions above, we obviously have “as λ → 0”.

C. Fixing phase and normalisation of perturbed eigenvectors

Eigenvectors in a complex vector space are only defined up to a complex scalar or, alterna-
tively, up to phase and magnitude. Hence, we impose the following conditions relating the
perturbed eigenvalues and eigenvectors to the unperturbed ones.
We require, for all λ ∈ (−ε, ε), all n and all δ,

(i) Imhenδ |enδ (λ)i = 0

(ii) kenδ (λ)k2 = 1.

Inserting the formal power series ansatz into these conditions yields
(k)
(i) Imhenδ |nδ i = 0 for k = 1, 2, . . .
(1) (2) (1)
(ii) 0 = 2λ Rehenδ |nδ i + λ2 2 Rehenδ |nδ i + knδ k2 + O(λ3 ).
13
German for “educated guess”.

– 91 –
Since (ii) holds for all λ ∈ (−ε, ε), we must have
(1) (2) (1)
Rehenδ |nδ i = 0, 2 Rehenδ |nδ i + knδ k2 = 0.
(1) (2)
Since we know from (i) that Imhenδ |nδ i = 0 and Imhenδ |nδ i = 0, we can conclude
(1) (2) (1)
henδ |nδ i = 0, henδ |nδ i = − 12 knδ k2 .
(1)
That is, nδ is orthogonal to enδ and
(2) (1)
nδ = − 21 knδ k2 enδ + e

for some e ∈ span({enδ })⊥ .

D. Order-by-order decomposition of the perturbed eigenvalue problem

Let us insert our formal power series ansatz into the perturbed eigenvalue equation. On
the left-hand side, we find
(1) (2)
Hλ enδ (λ) = (H0 + λW )(enδ + λnδ + λ2 nδ + O(λ3 ))
(1) (1) (2)
= H0 enδ + λ(W enδ + H0 nδ ) + λ2 (W nδ + H0 nδ ) + O(λ3 ),

while, on the right-hand side, we have

(1) (2) (1) (2)
hnδ (λ)enδ (λ) = (hn + λθnδ + λ2 θnδ + O(λ3 ))(enδ + λnδ + λ2 nδ + O(λ3 ))
(1) (1) (2) (1) (1) (2)
= hn enδ + λ(hn nδ + θnδ enδ ) + λ2 (hn nδ + θnδ nδ + θnδ enδ ) + O(λ3 ).

Comparing terms order-by-order yields

(H0 − hn )enδ = 0
(1) (1)
(H0 − hn )nδ = −(W − θnδ )enδ
(2) (1) (1) (2)
(H0 − hn )nδ = −(W − θnδ )nδ + θnδ enδ .

Of course, one may continue this expansion up to the desired order. Note that the zeroth
order equation is just our unperturbed eigenvalue equation.

E. First-order correction
To extract information from the first-order equation, let us project both sides onto the
unperturbed eigenvectors enα (i.e. apply henα | · i to both sides). This yields
(1) (1)
henα |(H0 − hn )nδ i = −henα |(W − θnδ )enδ i.

By self-adjointness of H0 , we have
(1) (1) (1)
henα |(H0 − hn )nδ i = h(H0 − hn )∗ enα |nδ i = h(H0 − hn )enα |nδ i = 0.

Therefore,
(1) (1)
0 = −henα |W enδ i + henα |θnδ enδ i = −henα |W enδ i + θnδ δαδ

– 92 –
and thus, the first-order eigenvalue correction is
(1)
θnδ = henδ |W enδ i.

Note that the right-hand side of the first-order equation is now completely known and
(1)
hence, if H0 −hn were invertible, we could determine nδ immediately. However, this is only
possible if the unperturbed eigenvalue hn is non-degenerate. More generally, we proceed as
follows. Let E := EigH0 (hn ). Then, we can rewrite the right-hand side of the first-order
equation as
(1) (1)
−(W − θnδ )enδ = − idH (W − θnδ )enδ
(1)
= −(PE + PE ⊥ )(W − θnδ )enδ
d(n)
(1) (1)
X
=− henβ |(W − θnδ )enδ ienβ − PE ⊥ W enδ + θnδ PE ⊥ enδ
β=1
= −PE ⊥ W enδ
(1)
so that we have (H0 − hn )nδ ∈ E ⊥ . Note that the operator

PE ⊥ ◦ (H0 − hn ) : E ⊥ → E ⊥

is invertible. Hence, the equation

(1)
PE ⊥ (H0 − hn )PE ⊥ nδ = −PE ⊥ W enδ

is solved by
(1)
PE ⊥ nδ = −PE ⊥ (H0 − hn )−1 PE ⊥ W enδ .
(1)
The “full” eigenvector correction nδ is given by
d(n)
(1) (1)
X
idH nδ = (PE + PE ⊥ )nδ = cδβ enβ − PE ⊥ (H0 − hn )−1 PE ⊥ W enδ ,
β=1

where the coefficients cδβ cannot be fully determined at this order in the perturbation.
What we do know is that our previous fixing of the phase and normalisation of the perturbed
(1)
eigenvectors implies that nδ is orthogonal to enδ , and hence we must have cδδ = 0.

F. Second-order eigenvalue correction

Here we will content ourselves with calculating the second-order correction to the eigen-
values only, since that is the physically interesting formula. As before, we proceed by
projecting both sides of the second-order equation onto an unperturbed eigenvector, this
time specifically enδ . We find

(2) (1) (1) (1) (2)

– 93 –
(1)
and recalling that henδ |nδ i = 0 and henδ |enδ i = 1, we have
(2) (1)
θnδ = henδ |W nδ i.
(1)
Plugging in our previous expression for nδ yields
d(n)
(2)
X
−1

θnδ = enδ W
cδβ enβ − W PE ⊥ (H0 − hn ) PE ⊥ W enδ
β=1
d(n)
X
= cδβ henδ |W enβ i − henδ |W PE ⊥ (H0 − hn )−1 PE ⊥ W enδ i
β=1
d(n)
(1)
X
= cδβ θnβ δδβ − henδ |W PE ⊥ (H0 − hn )−1 PE ⊥ W enδ i
β=1
= −henδ |W PE ⊥ (H0 − hn )−1 PE ⊥ W enδ i

since cδδ = 0. One can show that the eigenvectors of H0 (or any other self-adjoint operator)
form an orthonormal basis of H. In particular, this implies than we can decompose the
identity operator on H as
∞ d(n)
X X
idH = henβ | · ienβ .
n=1 β=1

(2)
By inserting this appropriately into our previous expression for θnδ , we obtain

∞ d(m)
X |hemβ |W enδ i|
(2)
X
θnδ =− .
hm − hn
m=1 β=1
m6=n

Putting everything together, we have the following second-order expansion of the perturbed
eigenvalues
(1) (2)
hnδ (λ) = hn + λθnδ + λ2 θnδ + O(λ3 )
∞ d(m)
X X |hemβ |W enδ i|
2
= hn + λhenδ |W enδ i − λ + O(λ3 ).
hm − hn
m=1 β=1
m6=n

Remark 8.10 . Note that, while the first-order correction to the perturbed nδ eigenvalue
only depends on the unperturbed nδ eigenvalue and eigenvector, the second-order correction
draws information from all the unperturbed eigenvalues and eigenvectors. Hence, if we try
to approximate a relativistic system as a perturbation of a non-relativistic system, then the
second-order corrections may be unreliable.

– 94 –
9 Case study: momentum operator

We will now put the machinery developed so far to work by considering the so-called
momentum operator for two cases: a compact interval [a, b] ⊂ R, and a circle. As the
name suggests, this operator is meant to be the QM observable who’s eigenvalues are the
momenta of the system. It is clear, therefore, that we require (recall H = L2 (Rd ) up to
unitary equivalence)
P : DP → L2 (Rd )
to be self adjoint.
We will specialise to the case of d = 1 in order to simplify things, while also demon-
strating the main ideas. The concepts can be extended to higher values of d. We will also
set ~ = 1 throughout this section.

9.1 The Momentum Operator

Definition. The momentum operator is an operator P given by

P : DP → L2 (R)
ψ 7→ (−i)ψ 0 ,

where the prime indicates a derivative.

The first obvious question is why does this deserve its name? (i.e. how is it related
to what we know classically as the momentum?) The answer, unfortunately, can not yet
be provided in full detail as it requires us to know the spectral theorem and Stone-von
Neumann theorem. However these details are provided here as they will help later when
discussing these theorems. For now we must just take it in faith.14
There is yet another important question we must ask: how do we choose DP ? The
immediate response might be ‘such that the derivative is square integrable.’ However, this
is not good enough. We also require that P be self adjoint and, as we have seen previously,
the concept of self adjointness depends heavily on the domains considered.
Luckily, not all hope is lost. The method will be as follows: guess a reasonable DP and
then search for a self adjoint extension, should one exist. Before doing so, though, we will
first introduce some new definitions that will prove invaluable.

9.2 Absolutely Continuous Fucntions and Sobolev Spaces

During the calculations that follow we will naturally encounter three spaces:

(i) The space of once-continuously differential functions over some interval, I; C 1 (I),

(ii) The Sobolev space H 1 (I), and

(iii) The space of absolutely continuous functions; AC(I).

14
Very unlike Dr. Schuller.

– 95 –
As we shall see they are related via

C 1 (I) ⊆ H 1 (I) ⊆ AC(I),

and so they will provide a convenient way to compare the domains DP , DP ∗ , etc, to test
for self adjointness.

Definition. Let I ⊂ R. A function ψ : I → C is called absolutely continuous if there exists

a Lebesgue integrable function ρ : I → C such that
Z x
ψ(x) = ψ(a) + ρ(y)dy,
a

for all compact subsets [a, x] ⊆ I.

Corollary 9.1. Given a absolutely continuous function, it is clear that ρ =a.e. ψ 0 , where the
almost everywhere condition comes from the fact that a Lebesgue integral does not distinguish
two elements that differ by a measure zero.

Definition. The set of absolutely continuous functions is simply

AC(I) := {ψ ∈ L2 (I) | ψ is absolutely continuous}.

Definition. Let Ω ⊂ R be open, and ψ : Ω → C be Lebesgue measurable. ψ is called

p-locally integrable if, for 1 ≤ p ≤ ∞,
Z
|ψ(x)|p dx < ∞,
K

for all compact subsets K ⊂ Ω. The set of all functions is

Lploc (Ω) := {ψ : Ω → C | ψ measurable, ψ|K ∈ Lp (K), ∀K ⊂ Ω, K compact}.

Remark 9.2 . For case p = 1, we just call ψ locally integrable.

Theorem 9.3. Every ψ ∈ Lp (Ω) for 1 ≤ p ≤ ∞ is locally integrable. In other words

Lp (Ω) ⊂ L1loc (Ω).

Definition. A function ψ ∈ L1loc (Ω) is called weakly differentiable if there exists a ρ ∈

L1loc (Ω) such that Z Z
ψ(x)ϕ0 (x)dx = − ρ(x)ϕdx,
Ω Ω
for all ϕ ∈ Cc∞ (Ω)15 . This function is known as the weak derivative of ψ and is denoted by
ρ := ψ 0 .

Corollary 9.4. Note that for any weakly differentiable function the integration by parts
result, Z Z
ψ(x)ϕ0 (x)dx = − ψ 0 (x)ϕ(x)dx,
Ω Ω
holds for all ϕ ∈ Cc∞ (Ω).
15
The subscript indicates that ϕ vanishes at the limits of integration.

– 96 –
Corollary 9.5. Given that ψ ∈ C k (Ω), ϕ ∈ Cc∞ (Ω), we can show by induction that
Z Z
(α) α
ψ(x)ϕ (x) = (−1) ψ (α) (x)ϕ(x),
Ω Ω

where ψ (α) (x) means the α-order derivative of ψ(x).

Remark 9.6 . In the above Corollary we have used the fact that we are only considering
one dimensional problems here. The expression is much the same for higher dimensional
problems, however one has to take into account the different derivative directions.

Definition. Let Ω ⊂ R be open, k ∈ N and 1 ≤ p ≤ ∞. The Sobolev Space is the space

with set
W k,p := {ψ ∈ Lp (Ω) ∩ W k (Ω) | ψ (α) ∈ Lp (Ω), ∀|α| ≤ k},
where W k is the set of all locally integrable functions that also have weak derivates of order
α for all |α| ≤ k. We introduce the notation H k (Ω) := W k,2 (Ω).

Remark 9.7 . Sobolev spaces can be made into Banach spaces by equipping them with a
norm and H k (Ω) can be made into a Hilbert space.

Proposition 9.8. We can rewrite the space H 1 (Ω) as

H 1 (Ω) := {ψ ∈ AC(Ω) | ψ 0 ∈ L2 (Ω)},

where ψ 0 denotes the normal notion of derivative.

Proof. See Theorem 7.13 in ‘A first Course in Sobolev Spaces, Giovanni Leoni’.

9.3 Momentum Operator on a Compact Interval

Let’s now consider the case where the physical space is some compact interval in R (i.e.
we have a particle moving along the bottom of a well). W.l.o.g. take [0, 2π] =: I, the
justification for which will is clear from the fact that next we will consider a circle.
We now need to come up with a reasonable guess for the domain DP . First recall

P : DP → L2 (I)
ψ 7→ (−i)ψ 0 .

It is therefore reasonable to restrict ourselves to ψ ∈ C 1 (I). Equally, physically we expect

the function to vanish at the boundaries (the walls of the well). This gives us our first guess

DP := {ψ ∈ C 1 (I) | ψ(0) = 0 = ψ(2π)} =: Cc1 (I).

The question still remains, though, as to whether P is self adjoint. Recalling the results
of Lecture 7, it is first instructive to see if P is symmetric.

– 97 –
A. Symmetric?
Let ψ, ϕ ∈ DP , then
Z 2π
hψ, Pϕi = dxψ(x)(−i)ϕ0 (x)
0
Z2π 2π
dx(−i)ψ 0 (x)ϕ(x) − i ψ(x)ϕ(x) 0

=−
Z 2π0
= dx(−i)ψ 0 (x)ϕ(x)
0
= hPψ, ϕi,

where integration by parts was used. So, yes P is symmetric.

B. Self Adjoint?
From above we know that P ⊆ P ∗ , and so DP ⊆ DP ∗ , so we need to ask the question of how
P ∗ behaves outside the domain DP . The obvious answer is to just extend the definition to
be

P ∗ : DP ∗ → L2 (I)
ψ 7→ (−i)ψ 0 .

Note that the ψ here is not necessarily the same as the ψ in above. The same symbol is
just used and the context tells us where it lives.
All that is left to check is the domain DP ∗ . From the definition of the adjoint we have

ψ ∈ DP ∗ =⇒ ∃η ∈ L2 (I) : ∀ϕ ∈ DP : hψ, Pϕi = hη, ϕi,

with η := P ∗ ψ. Before proceeding further with the calculation first introduce a function
N : I → C such that η =a.e. N 0 . Note N is Lebesgue integrable and that the almost
everywhere condition is sufficient as η appears in a Lebesgue integral. Therefore we have,
Z 2π Z 2π
0
dxψ(x)(−i)ϕ (x) = dxN 0 (x)ϕ(x)
0 0
Z 2π
2π
dx ψ(x)(−i)ϕ0 (x) + N (x)ϕ0 (x) = N (x)ϕ(x) 0

0
Z 2π
dx ψ(x) − iN (x) ϕ0 (x) = 0

−i
0
hψ − iN, ϕ0 i = 0,

which tells us that

ψ − iN ∈ {ϕ0 | ϕ ∈ DP }⊥ .
This does not appear to have got us any closer to determining the domain DP ∗ . However
consider the following two Lemmas.
R 2π
Lemma 9.9. {ϕ0 | ϕ ∈ DP } = {ξ ∈ C 0 (I) | 0 ξ(x)dx = 0}.

– 98 –
R 2π
Proof. Let A := {ϕ0 | ϕ ∈ DP } and B := {ξ ∈ C 0 (I) | 0 ξ(x)dx = 0}.
Now consider a ϕ0 ∈ A, then
Z 2π
2π
ϕ0 (x)dx = ϕ(x) 0 = 0,

0

so clearly ξ := ϕ and A ⊆ B.
Now consider a ξ ∈ B and define
Z x
ϕ(x) := ξ(y)dy.
0

Then, since ξ ∈ C 0 (I) it follows that ϕ ∈ C 1 (I). It also follows that ϕ(0) = 0 = ϕ(2π) and
so ϕ0 ∈ A and B ⊆ A.

Lemma 9.10. Let {1} denote the set consisting of the element 1 ∈ L2 (I) with 1(x) = 1C
for all x ∈ I. Then
{ϕ0 | ϕ ∈ DP } = {1}⊥ .
Proof. From the previous Lemma we have

{ϕ0 | ϕ ∈ DP } = {ξ ∈ C 0 (I) | h1, ξi = 0}

= {ξ ∈ C 0 (I) | h1, ξi = 0}
= {ξ ∈ L2 (I) | h1, ξi = 0}
= {1}⊥ ,

where the fact that C 0 (I) is dense in L2 (I) to go from the second to third line.

Putting this all together we have

ψ − iN ∈ A⊥ = A⊥
= (A⊥ )⊥⊥
= (A⊥⊥ )⊥
⊥
=A
= {1}⊥⊥
= {1}
= {C : I → C | x 7→ CC },

where CC is a constant in C. Recalling that N is Lebesgue integrable we see that

ψ(x) = CC + iN (x) ∈ AC(I),

and so
DP ∗ ⊆ AC(I).
Now recalling

P ∗ : DP ∗ → L2 (I)
ψ 7→ (−i)ψ 0 ,

– 99 –
and using Proposition 9.8, we have

DP ∗ ⊆ H 1 (I).

Finally we see that because all of the integration by parts results above were of the form
Z Z
dxψ 0 (x)ϕ(x) = ψ(x)ϕ0 (x),
I I

for arbitrary ϕ ∈ Cc1 (I).Now since ⊂ Cc1 (I), the integrals also hold for any ϕ ∈
Cc∞ (I)
Cc∞ (I), but this is just the condition for weak derivative and so we see that

DP ∗ = H 1 (I),

and
DP $ DP ∗ ,
so P is not self adjoint.

C. Essentially Self Adjoint?

We have managed to show that our initial guess for P is not self adjoint. The next step
is to ask if there is a self adjoint extension and if this extension is unique. Recall that a
symmetric operator has a unique self adjoint extension if it is essentially self adjoint (i.e.
P := P ∗∗ is self adjoint). We follow the same method as above,

ψ ∈ DP =⇒ ∀ϕ ∈ DP ∗ : hψ, P ∗ ϕi = hPψ, ϕi.

Now, recall that for a symmetric operator P ⊆ P ⊆ P ∗ so it’s clear that

hPψ, ϕi = hP ∗ ψ, ϕi
in the above. Writing as integrals we have

Z 2π Z 2π
0
dxψ(x)ϕ (x) = dx(−i)ψ 0 (x)ϕ(x)
0 0
Z 2π 2π
dx ψ(x)ϕ0 (x) − ψ(x)ϕ0 (x) = i ψ(x)ϕ(x) 0

−i
0
0 = ψ(2π)ϕ(2π) − ψ(0)ϕ(0),

where again integration by parts has been used. We need to be careful in what conclusions
we draw from this final statement, though. ϕ ∈ DP ∗ = H 1 (I), which places no restrictions
on the values of ϕ on the boundary, nor does it make any conditions between the two values
ϕ(0) and ϕ(2π) — they are independently arbitrary. We must, therefore, conclude that

ψ(2π) = ψ(2π) = 0 = ψ(0) = ψ(0),

and so, at best,

DP = {ψ ∈ H 1 (I) | ψ(2π) = 0 = ψ(0)} $ DP ∗ ,
∗
and so P =6 P ∗ , which, after taking the adjoint of both sides, tells us that P =
6 P , and so
P is not even essentially self adjoint.

– 100 –
D. Defect Indices
We only have one tool left to check for a self adjoint extension of P, check the defect indices
to see if a self adjoint extension even exists. Recall

d+ := dim ker(P ∗ − i) , d− := dim ker(P ∗ + i) ,

and a symmetric operator has a (not necessarily unique) self adjoint extension if d+ = d− .
We therefore need to determine how many ψ ∈ DP ∗ lie in ker(P ∗ ∓ i):

(P ∗ ∓ i)ψ = 0
−iψ 0 ∓ iψ = 0
ψ(x) = a± e∓x

for a+ , a− ∈ C. There is only one solution for each and so d+ = 1 = d− . We therefore know
that there does exist at least one self adjoint extension of P 16 , however we don’t know the
form of any of them.17

Remark 9.11 . If instead of a compact interval we take a half line I = [a, ∞), then d+ 6= d−
and so there is no self adjoint extension of P, meaning there is no notion of a QM momentum
in this case. Note however that people often talk about free particles along an infinite line
in QM, however they always require the wave function (ψ) to vanish at ±∞. This is clearly
just the same as taking a large, yet finite, compact interval I = [a, b].

9.4 Momentum Operator on a Circle

We now want to repeat all of the above but for a circle instead of a finite line segment.
Fortunately almost all the work is done, the only slight difference is in the definition of DP :

DP := {ψ ∈ C 1 (I) | ψ(2π) = ψ(0)},

which is exactly the same as before apart from now we do not require ψ to vanish at the
boundary. We still have

P : DP → L2 (R)
ψ 7→ (−i)ψ 0 ,

so it follows that PI $ Pc , where the I and c denote interval and circle respectively. In
other words Pc is an extension of PI .

A. Symmetric?
Repeating the steps from above it is clear that P is still symmetric. Note however, it is
symmetric for a different reason: before we had ψ(x)ϕ(x) = 0 as both ψ and ϕ vanished

at the limits, whereas now it holds simply because ψ(2π) = ψ(0) and likewise for ϕ.
16
Whew!
17
Not whew!

– 101 –
B. Self Adjoint?
As before we have

ψ ∈ DPc∗ =⇒ ∀ϕ ∈ DPc : hψ, Pc ϕi = hPc∗ ψ, ϕi.

However, recalling that the adjoint flips the inequality sign we have Pc∗ ⊆ PI∗ and therefore
DPc∗ ⊆ H 1 (I). We can, therefore, replace the unknown Pc∗ with the known PI∗ in the final
part of the above, i.e.
hPc∗ ψ, ϕi = hPI∗ ψ, ϕi.
Then, following the exactly as before, we arrive at
2π
0 = i ψϕ(x) 0

= ψ(2π) − ψ(0) ϕ(0)
=⇒ ψ(2π) = ψ(0),

giving us18
DPc∗ := {ψ ∈ H 1 (I) | ψ(2π) = ψ(0)} =: Hcyc
1
(I),
so DPc $ DP
∗ , and therefore P is not self adjoint.
c c

C. Essentially Self Adjoint?

Again, as before, we have Pc ⊆ Pc ⊆ Pc∗ and

ψ ∈ DPc =⇒ ∀ϕ ∈ DPc∗ : hψ, Pc∗ ϕi = hPc ψ, ϕi = hPc∗ ψ, ϕi,

which results in
2π
0 = i ψϕ(x) 0

= ψ(2π) − ψ(0) ϕ(0)
=⇒ ψ(2π) = ψ(0),

and so
1
DPc := Hcyc (I) = DPc∗ ,
so we conclude that Pc is essentially self adjoint and Pc is the unique self adjoint extension.
To summarise, we have found the momentum operator on a circle:
1
PS 1 : Hcyc (I) → L2 (R)
ψ 7→ (−i)ψ 0 .

18
Note we are OK extend the domain to all of H 1 (I) provided we impose the conditions above.

– 102 –
10 Inverse spectral theorem

This section is devoted to the development of all the notions and results necessary to
understand and prove the spectral theorem, stated below.

Theorem 10.1 (Spectral theorem). For every self-adjoint operator A : DA → H there is a

unique projection-valued measure PA : σ(OR ) → L(H) such that
Z Z
A = iR dPA ≡ λ PA (dλ),
R R

where iR : R ,→ C is the inclusion of R into C.

While useful in theory, existence results are often of limited use in practice since they
usually only tell us that something exists, and not how to construct it. However, we should
note here that the proof of the spectral theorem is, in fact, constructive in nature. Hence,
given any self-adjoint operator A, we will be able to explicitly determine its associated
projection-valued measure PA along the following steps.

ψ : σ(OR ) → R given by
(i) For each ψ ∈ H, construct the real-valued Borel measure µA
Z λ+δ
µA
ψ ((−∞, λ]) := lim lim dt Imhψ|RA (t + iε)ψi,
δ→0+ ε→0+ −∞

where RA : ρ(A) → L(H) is the resolvent map of A. This is know as the Stieltjes
inversion formula. Note that while not every element in σ(OR ) is of the form (−∞, λ],
such Borel measurable sets do generate the entire σ(OR ) via unions, intersections and
set differences. Hence, the value of µAψ (Ω) for Ω ∈ σ(OR ) can be determined by
applying the corresponding formulae for measures, namely σ-additivity, continuity
from above and measure of set differences.

ψ,ϕ : σ(OR ) → C by
(ii) For all ψ, ϕ ∈ H, define the complex-valued Borel measure µA

µA 1 A A A A
ψ,ϕ (Ω) := 4 (µψ+ϕ (Ω) − µψ−ϕ (Ω) + iµψ−iϕ (Ω) − iµψ+iϕ (Ω)).

(iii) Define the projection-valued measure PA : σ(OR ) → L(H) by requiring PA (Ω), for
each Ω ∈ σ(OR ), to be the unique map in L(H) satisfying
Z
∀ ψ, ϕ ∈ H : hψ|PA (Ω)ϕi = χΩ dµA ψ,ϕ .
R

We will now make all the notions and constructions used herein precise. In fact, we
will present the relevant definitions and results by taking the inverse route, starting with
projection-valued measures and arriving at their associated self-adjoint operators, obtaining
(and proving) what we will call the inverse spectral theorem.

– 103 –
10.1 Projection-valued measures
Projection-valued measures are, unsurprisingly, objects sharing characteristics of both mea-
sures and projection operators.

Definition. A map P : σ(OR ) → L(H) is called a projection-valued measure if it satisfies

the following properties.

(i) ∀ Ω ∈ σ(OR ) : P(Ω)∗ = P(Ω)

(ii) ∀ Ω ∈ σ(OR ) : P(Ω) ◦ P(Ω) = P(Ω)

(iii) P(R) = idH

(iv) For any pairwise disjoint sequence {Ωn }n∈N in σ(OR ) and any ψ ∈ H,
∞
X ∞
[
P(Ωn )ψ = P Ωn ψ.
n=0 n=0

Remark 10.2 . Note in the final condition we included ψ ∈ H as, for the case countably
infinite n ∈ N we need to check convergence, which involves using the norm. Without the ψ
we would need to use the norm on L(H) which may prove difficult. However, by including
the ψ we can work with the norm on H itself.

Lemma 10.3. Let P : σ(OR ) → L(H) be a projection-valued measure. Then, for any
Ω, Ω1 , Ω2 ∈ σ(OR ),

(i) P(∅) = 0, where by 0 we mean 0 ∈ L(H)

(ii) P(R \ Ω) = idH − P(Ω)

(iii) P(Ω1 ∪ Ω2 ) = P(Ω1 ) + P(Ω2 ) − P(Ω1 ∩ Ω2 )

(iv) P(Ω1 ∩ Ω2 ) = P(Ω1 ) ◦ P(Ω2 )

(v) if Ω1 ⊆ Ω2 , then ran(P(Ω1 )) ⊆ ran(P(Ω2 )).

Proof. Let Ω, Ω1 , Ω2 ∈ σ(OR ),

(i)

P (∅)ψ = P (∅ ∪ ∅)ψ = P (∅)+P (∅) ψ = 2P (∅)ψ
∴ P (∅)ψ = 0H =⇒ P (∅) = 0.

(ii)

P (R) = P (R \ Ω) ∪ Ω = P (R \ Ω) + P (Ω)
∴ P (R \ Ω) = idH −P (Ω).

where we used the fact that (R \ Ω) ∩ Ω = ∅.

– 104 –
(iii)

P (Ω1 ) = P (Ω1 ∩ Ω2 ) ∪ (Ω1 \ Ω2 ) = P (Ω1 ∩ Ω2 ) + P (Ω1 \ Ω2 ) ,
and similarly for P (Ω2 ). Also

P (Ω1 \ Ω2 ) + P (Ω2 \ Ω1 ) + P (Ω1 ∩ Ω2 ) = P (Ω1 \ Ω2 ) ∪ (Ω2 \ Ω1 ) ∪ (Ω1 ∩ Ω2 )
= P (Ω1 ∪ Ω2 ).

Putting this all together gives the result.

(iv) First consider Ω1 ∩ Ω2 = ∅. Then, using (ii) from the definition we have

[P (Ω1 ∪ Ω2 )]2 = [P (Ω1 ) + P (Ω2 )]2

P (Ω1 ∪ Ω2 ) = P (Ω1 ) + P (Ω2 ) + P (Ω1 ) ◦ P (Ω2 ) + P (Ω2 ) ◦ P (Ω1 )
∴ P (Ω1 ) ◦ P (Ω2 ) = −P (Ω2 ) ◦ P (Ω1 )
=⇒ P (Ω1 ) ◦ P (Ω2 ) ◦ P (Ω2 ) = −P (Ω2 ) ◦ P (Ω1 ) ◦ P (Ω2 )
P (Ω1 ) ◦ P (Ω2 ) = P (Ω2 ) ◦ P (Ω2 ) ◦ P (Ω1 )
P (Ω1 ) ◦ P (Ω2 ) = P (Ω2 ) ◦ P (Ω1 )
∴ P (Ω1 ) ◦ P (Ω2 ) = 0 ∀Ω1 ∩ Ω2 = ∅.

Now from P (Ω1 ) = P (Ω1 \ Ω2 ) ∪ (Ω1 ∩ Ω2 ) = P (Ω1 \ Ω2 ) + P (Ω1 ∩ Ω2 ), we have

P (Ω1 ) ◦ P (Ω2 ) = [P (Ω1 \ Ω2 ) + P (Ω1 ∩ Ω2 ) ] ◦ [P (Ω2 \ Ω1 ) + P (Ω1 ∩ Ω2 ) ]
= P (Ω1 ∩ Ω2 ) ◦ P (Ω1 ∩ Ω2 )
= P (Ω1 ∩ Ω2 )

where we have made use of the fact that (Ω1 \ Ω2 ) ∩ (Ω2 \ Ω1 ) = ∅ etc.

(v) If Ω1 ⊂ Ω2 we have

P (Ω2 ) = P (Ω2 \ Ω1 ) ∪ Ω1 ) = P (Ω2 \ Ω1 ) + P (Ω1 ),

which along with the fact that P (Ω) ≥ 0 gives the result.

Note most of these properties make sense simply by thinking of P (Ω) as the area of
the set Ω ∈ σ(OR ). For example P (Ω1 ) = P (Ω1 \ Ω2 ) + P (Ω1 ∩ Ω2 ) is

Ω1 Ω2 Ω1 Ω2

A\B A∩B B\A A\B A∩B B\A

– 105 –
Remark 10.4 . As noted before, it suffices to know P (−∞, λ] for all λ ∈ R, and for this

reason a new notation is introduced;

P (λ) := P (−∞, λ] .
However, as they are written they appear to have different domains. We therefore give a
new name to P (λ); the resolution of the identity.

10.2 Real and Complex Valued Borel Measures Induced by a PVM

Definition. For all ψ, ϕ ∈ H we define the C-valued measure
µψ,ϕ : σ(OR ) → C
Ω 7→ µψ,ϕ (Ω) := hψ, P (Ω)ϕi.
Definition. For all ψ ∈ H we define the R-valued measure
µψ := µψ,ψ
Proof. that µψ (Ω) ∈ R.
µψ (Ω) = hψ, P (Ω)ψi
= hP (Ω)ψ, ψi
= hψ, P (Ω)ψi
= µψ (Ω),
where we have used the fact that P is self adjoint to go from the first to the second line.

10.3 Integration With Respect to a PVM

We now wish to make sense of the operator R f dP for measurable f : R → C. We will
R

build this up in three steps:

(i) For simple f ,

(ii) For bounded f , and

(iii) For not necessarily bounded19 f .

Remark 10.5 . As we shall see, if f is bounded (in the sense that there exists a a ∈ R such
that |f (x)| < a for all x ∈ R) that the integral part of the operator will have nice properties.
For example it will be linear in f , i.e.
Z Z Z
(αf + g)dP = α f dP + gdP,
R R R
for α ∈ C. However, if f is unbounded, domain issues will destroy the equality above. It is
important to note, though, that it is exactly the latter case we need as
f = iR : R ,→ C
x 7→ x
in the Spectral theorem, which is clearly unbounded.
19
Recall footnote 5 from lecture 2: to us the term ‘unbounded’ means definitely not bounded.

– 106 –
A. Simple f
Recall a simple function is a measurable function that takes a finite number of results in
the target, i.e. if f : R → C then

f (R) = {f1 , ..., fN } ⊂ C

for some N ∈ N. This allows us to rewrite f as

N
X
f= fn χΩn ,
n=1

where Ωn := preimf ({fn }) and χ is the characteristic function.

Definition. For simple f : R → C and PVM P we define

Z N
X
f dP := fn P (Ωn ).
R n=1
R
Proposition 10.6. For simple f , R f dP is linear in f .

Proof. Let S(R, C) denote the set of all simple functions f : R → C. We can make this set
into a C-vector space by inheriting the addition and s-multiplication from C, namely define

(f + g)(x) := f (x) + g(x), (α · f )(x) := α · f (x),

for all f, g ∈ S(R, C), α ∈ C.

Now consider the preimage part: As f and g are both simple it follows that

(f + g)(x) = fn + gn

for some fn , gn ∈ C, so the preimage term becomes

preimf +g {fn + gn } = {x} = preimf {fn }

= preimg {gn }.

It follows trivially, then, that

Z Z Z
(f + g)dP = f dP + gdP.
R R R

A similar method gives the α ∈ C condition of linearity.

Remark 10.7 . Observe that χΩ for any Ω ∈ σ(OR ) is simple (it only takes the values 0 or
1), and hence Z
χΩ dP = 1 · P (Ω) + 0 · P (σ \ Ω) = P (Ω).
R

– 107 –
Remark 10.8 . Observe also that for any ψ, ϕ ∈ H,
N
* Z + D X E
ψ, f dP ϕ = ψ, fn P (Ωn )ϕ
R n=1
N
X
= fn hψ, P (Ωn )ϕi
n=1
XN
= fn µψ,ϕ (Ωn )
n=1
Z
=: f dµψ,ϕ ,
R

where Proposition 6.7 was used.

Definition. For simple f we can define the map

Z
dP : S(R, C) → L(H)
R Z
f 7→ f dP,
R

which, if we equip S(R, C) with the suppremum norm and L(H) with its operator norm,
has operator norm k R dP k = 1.
R

Proof. We have already shown that R f dP ∈ L(H) (i.e. it is linear), so we just need to
R

show the norm condition. First, let f ∈ S(R, C) and ψ ∈ H, then

Z 2 * Z Z +

f dP ψ = f dP ψ, f dP ψ
R H R R
N
DX N
X E
= fn P (Ωn )ψ, fm P (Ωm )ψ
n=1 m=1
D N
X E
= ψ, fn fm P (Ωn )P (Ωm )ψ
n,m=1
D N
X E
= ψ, fn fm δnm P (Ωm )ψ
n,m=1
N
X
= |f |2 hψ, P (Ωn )ψi
n=1
XN
= |f |2 µψ (Ωn )
n=1
Z
=: |f |2 dµψ
Z R

=⇒
f dP ψ
≤ kf k∞ kψkH ,
R H

– 108 –
where we have used the definition of the norm in terms of the inner product, the fact
that P is self adjoint, the fact that P (Ωn )P (Ωm ) = δnm for pairwise disjoint Ωn /Ωm , and
Proposition 6.27 along with the fact that kf k∞ := supx∈R f (x). The equality in the last
line can be assumed provided f and ψ are sufficiently chosen.
Thus we have
R
Z f dP

dP := sup R L(H)

R

f ∈S(R,C) kf k∞
R
R f dP ψ H

:= sup sup
f ∈S(R,C) ψ∈H kf k∞ kψkH
= 1.

B. Bounded Borel Functions

Definition. The set of all bounded, measurable functions is denoted

B(R, C) := {f : R → C | measurable, kf k∞ < ∞}.

Proposition 10.9. The set B can be made into a Banach space by defining the norm

kf kB := sup |f (x)|.
x∈R

Proof. We turn the set into a linear vector space in the usual manner; we inherit the addition
and s-multiplication from C.
Now prove kf kB is a norm. Comparing to the definition given at the bottom of Page
9, for f, g ∈ B(R, C) and z ∈ C we have

(i) Clearly kf kB ≥ 0.

(ii)

kf kB = 0
⇔ sup |f (x)| = 0
x∈R
⇔ f (x) = 0 ∀x ∈ R

so f = 0.

(iii)

kz · f kB := sup |z · f (x)|
x∈R
= |z| sup |f (x)|
x∈R
=: |Z| · kf kB .

– 109 –
(iv)

kf + gkB := sup |(f + g)(x)|

x∈R
= sup |f (x) + g(x)|
x∈R
≤ sup |f (x)| + sup |g(x)|
x∈R x∈R
=: kf kB + kgkB .

Now let {fn }n∈N be a Cauchy sequence in B(R, C), that is: ∀ > 0, ∃N ∈ N : ∀m, n ≥ N
we have
d(fn , fm ) := kfn − fm kB := sup |fn (x) − fm (x)| < .
x∈R
Now from the definition of the supremum we have

|fn (x) − fm (x)| ≤ kf − gkB ,

so it follows that the sequence {fn (x)}n∈N is a Cauchy sequence in C. But C is a complete
metric space so we know that this Cauchy sequence converges in C, i.e.

lim fn (x) = zx ∈ C.
n→∞

We can thus define a point-wise limit f of the sequence {fn }n∈N ⊂ B(R, C) as

f (x) := lim fn (x) = zx

x→∞

for all x ∈ R. Then, by equipping C and R with their respective Borel σ-algebras, Proposi-
tion 5.22 tells us that f is measurable.
Finally from the fact that B(R, C) ⊂ L(R, C), Theorem 2.8 tells us that f is bounded
and so B(R, C) is a Banach space.

Corollary 10.10. Observe that S(R, C) is in-fact a dense, linear subspace of B(R, C). Thus,
the BLT theorem tells us that we have a unique extension of the operator
Z
dP : S(R, C) → L(H)
R

to the domain B(R, C) with equal operator norm. That is, we have an operator

\
Z
dP : B(R, C) → L(H)
R
R
with \R dP = 1.

Corollary 10.11. By suitable definition we can turn the space B(R, C) into an C ∗ -algebra
and our operator then has the following properties

(i) Z
1dP = idH
R

– 110 –
(ii) Z Z Z
(f · g)dP = f dP ◦ gdP
R R R

(iii) Z Z ∗
f dP = f dP
R R

These properties collectively make the operator a C ∗ -algebra homomorphism.

C. General Borel Function

We now want to allow for the case that f is unbounded. We will write the following such
that it reduces to the above when f is bounded.

Definition. Let f : R → C be measurable, then we define the linear map

Z
f dP : DR f dP → H
R
R

where Z
2
DR := ψ ∈ H |f | dµψ < ∞ ⊆ H,

R f dP
R
is a dense, linear subspace. The linear map is defined via
Z Z
f dP ψ := lim fn dP ψ ,
R n→∞ R

where the sequence {fn }n∈N ⊆ B(R, C) defined by

fn := χ{x∈R|f (x)<n} f.

Remark 10.12 . Note that the map above includes the f , it is not just the integral defined
in the previous section. Note also for the case when f ∈ B(R, C) we just recover the case
above and we have DR f dP = H (i.e. we have L(H)). Otherwise it is a proper subset.
R

Remark 10.13 . The literature often introduces the notation

Df := DR f dP ,
R

however this could lead one to think of the domain of f itself, which here is R. We will
avoid this notation.

Remark 10.14 . The sequence {fn }n∈N can be thought of as ‘chopping’ f into bounded parts.

f
f2
f1

– 111 –
Lemma 10.15. The sequence {fn }n∈N is Cauchy in L2 (R). This in tern implies that the
R
sequence {( R fn dP )ψ}n∈N is Cauchy in H, which is required for the limit in the definition
to make sense, i.e. the result lies in H.

For this general case of a measurable f we have:

(i) As before Z Z ∗
f dP = f dP .
R R

(ii) For α ∈ C and f, g measurable,

Z Z Z
α f dP + gdP ⊆ (αf + g)dP,
R R R

where the equality holds only for bounded f . As explained earlier, the inequality
arises due to domain issues. We can now see this more explicitly from the defini-
tion of DR f dP ; just because f and g are both measurable, it does not mean that
R
their respective map domains will coincide. However, the domain for the LHS is
DR (|αf |+|g|)dP .
R

(iii) Z Z Z
f dP ◦ gdP ⊆ (f · g)dP,
R R R
again where the equality holds only when f and g are bounded.

10.4 The Inverse Spectral Theorem

We are now in a place where we can understand the inverse spectral theorem.

Definition. Given a PVM, P , we can construct a self adjoint operator AP as

Z
AP := idR dP,
R

where idR : R ,→ C is the inclusion map.

Proof.
Z
∗
(AP ) = idR dP
ZR
= idR dP
R
= AP ,

so AP is self adjoint.

– 112 –
11 Spectral theorem
The inverse spectral theorem tells us how to construct a self adjoint operator AP given a
projection valued measure P . The aim of this lecture is to do the opposite; given a self
adjoint operator A we want to find a PVM PA . We want the two methods to be in unison
— that is we want APA = A and PAP = P . We shall start by assuming that A can be
written in integral form, and then shall remove this restriction.

11.1 Measurable Function Applied To A Spectrally Decomposable Self Adjoint

Operator
The term spectrally decomposable means that is has integral form.

Definition. Let the self adjoint operator A be spectrally decomposable; i.e. there exists a
PVM P such that Z
A= idR dP,
R
then for any measurable function f : R → C, we define the operator

f (A) : DR f dP → H,

given by Z Z
f (A) := (f ◦ idR )dP ≡ f (λ)P (dλ).
R R
Remark 11.1 . The spectral theorem will show that every self adjoint operator A is spectrally
decomposable by virtue of a uniquely detemerined P .

Corollary 11.2. If f : R → R, then f (A) is again self adjoint.

Proof.
Z ∗
[f (A)]∗ := (f ◦ idR )dP
Z R

= (f ◦ idR )dP
ZR
= f dP
R
=: f (A).

Let us now consider two important examples.

Example 11.3 . Let A = R λP (dλ) be self adjoint. Then

Z
exp(A) := eλ P (dλ)
R

is self adjoint due to the previous Corollary. However

Z
exp(iA) := eiλ P (dλ)
R

– 113 –
is not self adjoint. This latter case is of high importance in QM, as can be seen by revisiting
Axiom 4 at the start.

Example 11.4 . Let A = R λP (dλ) be self adjoint and Ω ∈ σ(OR ). Then

Z
P (Ω) = χΩ dP
R

implies
Z
sq P (Ω) :=

(sq ◦ χΩ dP
ZR
≡ [χΩ (λ)]2 P (dλ)
ZR
= χΩ (λ)P (dλ)
R
= P (Ω),

which is one of the projection conditions.

11.2 Reconstruct PVM From a Spectrally Decomposable, Self Adjoint Opera-

tor
The key to reconstructing the associated PVM P is to consider the resolvents.

Definition. Given a spectrally decomposable operator A and an the resolvant set ρ(A),
we define
rz (A) = RZ (A) := (A − z idH )−1 ,
which is rewritten as

rz : R → C
1
λ 7→ ,
λ−z
and, due to the fact that A is spectrally decomposable, satisfies
Z Z
1
rz (A) = (rz ◦ idR )dP ≡ P (dλ).
R R λ−z

Note, using the results in the previous lecture, we have that for any ψ ∈ H,
Z
hψ, RA (z)ψi = ψ, (rz ◦ idR )dP ψ
Z R

= (rz ◦ idR )dµψ

ZR
1
≡ µψ (dλ).
R λ−z

Definition. A Herglotz function is an analytic complex function that maps the upper
half plane into itself, but need not be surjective or injective. They are also known as
Nevanlinna/Pick/R functions.

– 114 –
Theorem 11.5. The function

hψ, RA (·)ψi : C → C
Z
1
z 7→ µψ (dλ)
R λ−z
is Herglotz.

Proof. Recall
µψ : σ(OR ) → R+
0

is real-valued. Then using

1
Im (r) = (r − r),
2
we have
Z
1
Im hψ, RA (z)ψi = Im µψ (dλ)
R λ−z
Z
1 1 1
= − µψ (dλ)
2 R λ−z λ−z
z−z
Z
= µψ (dλ)
R 2|λ − z|2
Z
1
= Im(z) µψ (dλ).
R |λ − z|2
Then, since the fact that the integral is Lebesgue and the intergrand is non-negative and
so
Imhψ, RA (z)ψi > 0 ⇔ Im(z) > 0.

Recalling the start of last lecture, if we can find a way to construct µψ from our A then
we can use that to reconstruct P . The result of this is the previously mentioned Stieltjes
Inversion Formula, and it is obtained as follows.
Let t, ∈ R. Then, since A is self adjoint and so its spectrum is purely real, t+i ∈ ρ(A).
This allows us to act on it with RA . Thus, consider
1 t2 1 t2
Z Z Z
ε
lim dt Imhψ, RA (t + iε)ψi = lim dt µ (dλ)
2 ψ
ε→0 π t1 ε→0 π t1 R |λ − t − iε|
+ +
Z Z t2
1 ε
= lim dt µψ (dλ),
ε→0+ R π t1 (λ − t)2 + ε2

where Fubini’s Theorem20 has been used. The inner integral is a standard integral, with
result
1 t2 t − λ t2
Z
ε 1
dt = arctan .
π t1 (λ − t)2 + ε2 π ε t1
Now strictly, at this stage, we cannot simply pull the ε limit into this expression; we would
need to check that the above result is bounded first and then, by dominated convergence,
20
See Wiki

– 115 –
we can pull it in. This will turn out to be true, and so in order to simplify the following we
consider ε to be small here.
In order to work out the above expression, we can use the λ-graphs. Let’s plot both
terms (including the overall minus sign that comes with t1 ) on the same graph:

1
t1 −λ
− π1 arctan

2
ε

t1 t2 λ

1 t2 −λ

π arctan ε
− 12

Taking the limit and adding gives

1
2

0
t1 t2 λ
So we have Z t2
1 ε 1
lim dt 2 2
= χ(t1 ,t2 ) + χ[t1 ,t2 ] ,
ε→0 π
+
t1 (λ − t) + ε 2
and Z t2 Z
1 1
lim dt Imhψ, RA (t + iε)ψi = χ(t1 ,t2 ) + χ[t1 ,t2 ] µψ (dλ).
ε→0 π
+
t1 2 R

Finally, we have the Stieltjes Inversion Formula.

Theorem 11.6 (Stieltjes Inversion Formula). Given a spectrally decomposable, self adjoint
operator A and its associated resolvent map RA , we can construct a real-valued measure
Z λ+δ
1
µA

ψ (−∞, λ] = lim lim dt Imhψ, RA (t + iε)ψi.
δ→0 ε→0 π
+ +
−∞

– 116 –
Proof.
Z λ+δ Z
1 1
lim lim dt Imhψ, RA (t + iε)ψi = lim χ(−∞,λ+δ) + χ(−∞,δ] µψ (dλ)
δ→0 ε→0 π
+ +
−∞ δ→0 2 R
+
Z
= χ(−∞,λ] µψ (dλ)
R
= µψ (−∞, λ]),

where we used the fact that the χ(Ω) is bounded to move the limit inside the integral along
with the fact that
lim (−∞, λ + δ) = (−∞, λ].
δ→0+

Remark 11.7 . Note the fact that hψ, RA (t + iε)ψi is Herglotz with the fact that ε > 0 gives
us that µA
ψ ≥ 0, which is required for it to be a real-valued measure.

Remark 11.8 . If we already know that A is spectrally decomposable w.r.t. some PVM P ,
then we can recover P from A by virtue of: for any Ω ∈ σ(OR ) and for all ψ, ϕ ∈ H
Z
hψ, P (Ω)ϕi = χΩ dµψ,ϕ ,
R
where µψ,ϕ is obtained from µψ using the method given at the start of the previous lecture.

11.3 Construction Of PVM From A Self Adjoint Operator

We need to free ourselves from the fact that A is known to be spectrally decomposable
from the start. We could do this by trying to recreate the above method, i.e. arrive at the
Stieltjes Inversion Formula for an operator A, by showing that

(i) hψ, RA (·)ψi : C → C is Herglotz for any self adjoint A

(ii) hψ, P (Ω)ϕi := R χΩ dµψ,ϕ is indeed a PVM.

In order to prove these we first need a new theorem.

Theorem 11.9 (First-Resolvent Formula). For any operator A : DA → H and a, b ∈ ρ(A)
we have
RA (a) − RA (b) = (a − b)RA (a)RA (b) = (a − b)RA (b)RA (a).
Proof. Consider

RA (a) − (a − b)RA (a)RA (b) := (A − a)−1 − (a − b)(A − a)−1 (A − b)−1

= (A − a)−1 idH −(a − b)(A − b)−1

= (A − a)−1 idH −(a − A + A − b)(A − b)−1

= (A − a)−1 idH +(A − a)(A − b)−1 − (A − b)(A − b)−1

= (A − b)−1
= RA (b),

and similarly for the the other result.

– 117 –
The proof of (i) and (ii) above was on a problem sheet. Return and do this later.
Conclusion, the Spectral Theorem, Theorem 10.1, together with the recipe for the
construction of the PVM P from a self adjoint operator A holds.

11.4 Commuting Operators

The study of QM is teaming with so called commutators. However, they are not as simply
defined as is often erroneously assumed. In particular, the commutator between the position
and momentum given by
[Qi , Pj ] = i~δ i j
is not even defined, unless further provisions are given. This formula appears in the opening
sections of almost all QM textbooks though, and so it’s important we understand what is
meant.
One can happily write the commutator provided the operators involved are bounded
(which from the lecture 9 we see that at least Pj is not).

Definition. Let B1 , B2 ∈ L(H), i.e. they are bounded linear operators from H to H. Then
one may define
[B1 , B2 ] := B1 ◦ B2 − B2 ◦ B1 ,
where
[B1 , B2 ] ∈ L(H).

Remark 11.10 . The tuple (L(H), +, ·, [·, ·]) is a Lie algebra21 .

Corollary 11.11. Let A, B, C ∈ L(H) then the following holds

[A ◦ B, C] = [A, C] ◦ B + A ◦ [B, C].

Proof. Consider the action on some arbitrary ψ ∈ H.

[A ◦ B, C]ψ := (A ◦ B) ◦ (Cψ) − C ◦ (A ◦ Bψ)

= A ◦ (B ◦ Cψ) − C ◦ (A ◦ Bψ)
= A ◦ (C ◦ Bψ + [B, C]ψ) − C ◦ (A ◦ Bψ)
= (A ◦ C) ◦ Bψ + A ◦ [B, C]ψ − C ◦ (A ◦ Bψ)
= (C ◦ A + [A, C]) ◦ Bψ + A ◦ [B, C]ψ − C ◦ (A ◦ Bψ)
= [A, C] ◦ Bψ + A ◦ [B, C]ψ
= ([A, C] ◦ B + A ◦ [B, C])ψ,

where we used the associativity of the composition of maps. Then, as ψ was arbitrary, we
have our result.

Corollary 11.12. Let A and B be two operators. Then if one of the them is unbounded
the domain D[A,B] may only have a trivial definition, i.e. D[A,B] = {0H }.
21
See Dr. Schuller’s Lectures on the Geometric Anatomy of Theoretical Physics

– 118 –
Proof. Let A : DA → H be unbounded with DA ⊂ H, and define the bounded operator

Bϕ :H → H
α 7→ hϕ, αiψ =: `ϕ (α)ψ

for some fixed ϕ, ψ ∈ H where ψ ∈

/ DA . So we have ran(Bϕ ) = H \ DA . Then from the first
term in
[A, Bϕ ] := A ◦ Bϕ − Bϕ ◦ A,
it follows that
D[A,Bϕ ] = ran(Bϕ ) ∩ DA = {0H }.

Definition. Two bounded linear operators A, B ∈ L(H) are said to commute if

[A, B] = 0.

Corollary 11.13. Let A, B ∈ L(H) be commuting operators. Then if A is also non-

degenerate then any ψ that is an eigenvector of A is also an eigenvector of B. In other
words, the set of A’s eigenvectors is contained within the set of B’s.

Proof. Let ψ ∈ D \ {0} be an eigenvector of A with eigenvalue λ ∈ C. Then we have

[A, B]ψ := (A ◦ B − B ◦ A)ψ

= A(Bψ) − B(Aψ)
= A(Bψ) − λB(ψ),

where we have used the fact that B is linear. Then from the fact that [A, B] = 0 it follows
that Bψ is also an eigenvalue of A with eigenvalue λ. Finally from the fact that A is
non-degenerate it follows that Bψ = µψ must hold for some µ ∈ C.

However, as highlighted at the start of this section, we also want to look at situations
when one of the operators may not be bounded. In other words we want to know how to
extend the idea of commuting to

(i) A self adjoint and not necessarily bounded, B bounded.

(ii) Both A and B self adjoint and not necessarily bounded.

As is often the case in maths/physics problems, the strategy is to reduce the problem
to the known case. We then have three possible bounded, linear operators constructed from
A:

(i) From the Spectral Theorem we know that if A is self adjoint then there exists a unique
PVM P such that A is spectrally decomposable. Recall that, from Remark 10.7,
PA (Ω) ∈ L(H) for any Ω ∈ σ(OR ).

– 119 –
(ii) From the definition of the resolvent set, we have RA (z) ∈ L(H) for any z ∈ ρ(A).

(iii) Again, as A is self adjoint it is spectrally decomposable and so we can consider

exp(itA) for some t ∈ R, which was defined in Example 11.3. Again this is not a self
adjoint operator, but it is unitary, which means k exp(itA)k = 1.

Definition. Let A be self adjoint and B be bounded. A and B are said to commute if
either of the following holds

(i) [RA (z), B] = 0 for some z ∈ ρ(A).

(ii) [exp(itA), B] = 0 for some t ∈ R \ {0}.

Definition. Let A and B be self adjoint. They are said to commute is one of the following
holds

(i) [RA (zA ), RB (zB )] = 0 for some zA ∈ ρ(A) and zB ∈ ρ(B). This is known as the
Resolvent way.

(ii) [exp(itA), exp(isB)] = 0 for some t, s ∈ R \ {0}. This is known as the Weil way.

(iii) [PA (Ω), PB (Ω)] = 0 for all Ω ∈ σ(OR ). This is known as the Projector way.

Remark 11.14 . The literature normally uses a practical, yet misleading, notation at this
point. For any of the above we simply write [A, B] = 0 for commuting A and B. However
this commutator is not the same as the one defined at the start — i.e. it does not correspond
to A ◦ B − B ◦ A. Really we should write it slightly differently to highlight this, i.e. make
it red; [A, B].

Theorem 11.15. Let A and B be self adjoint and bounded. Then

[A, B] = 0 ⇔ [A, B] = 0.

– 120 –
12 Stone’s Theorem and Construction of Observables
In this lecture we aim to answer two questions by deriving and using Stone’s Theorem.
They are

(i) How arbitrary is the stipulation of Axiom 4; that the dynamics in the absence of a
measurement be controlled by

U (t) := exp(itH)

for some self adjoint operator H?

(ii) How does one practically construct observables, including the question of how to find
the correct domain such that the operator is at least essentially self adjoint?

Remark 12.1 . Clearly for (i) we want U (t) ◦ U (s) = U (t + s) and U (0) = idH .

12.1 One Parameter Groups and Their Generators

Definition. A group is the double (G, ♦), where G is a set and ♦ : G → G satisfying:

(i) For all g, h, k ∈ G, (Associativity)

(g♦h)♦k = g♦(h♦k).

(ii) There exists e ∈ G such that for all g ∈ G (Neutral Element)

g♦e = e♦g = g.

(iii) For all g ∈ G there exists g −1 ∈ G such that (Inverse)

g♦g −1 = g −1 ♦g = e.

Definition. A Abelian group is a group is one whose group operation is symmetric. That
is for all g, h ∈ G
g♦h = h♦g.

Remark 12.2 . Abelian groups are also known as commutative groups and the condition is
refered to as the commutativity of the elements with respect to the group operation.

Example 12.3 . The real numbers equipped with addition form an Abelian group, with e =
0 ∈ R and g −1 = −g.
Note it is important that we consider all of R, and not just the positive numbers, as in
the latter case the inverse would not lie in the group.

Example 12.4 . The set R \ {0} form an Abelian group with respect to multiplication, with
e = 1 and g −1 = 1/g.
Note here we have to exclude 0 as 1/0 is not an element of R.

– 121 –
Example 12.5 . The set R \ {0} is not a group with respect to division as it fails to satisfy
associativity.

Definition. A one-parameter group is a group whose underlying set is

G = {U (t) | t ∈ R},

and whose group operation satisfies

U (t)♦U (s) = U δ(t, s) ,

for some δ : R × R → R.

Remark 12.6 . Unless the group is Abelian then δ(s, t) 6= δ(t, s).

We will only deal with Abelian one-parameter groups, in which case one can always
reparameterise so that
U (t)♦U (s) = U (t + s)
where the commutativity with respect to ♦ is inherited from the commutativity with respect
to +. We also choose the parameterisation such that U (0) = e. In particular, we will look
at unitary one-parameter groups, i.e. those with

G = {U (t) ∈ L(H) | t ∈ R, U ∗ (t)U (t) = idH kU (t)k = 1},

that are strongly continuous in the parameter,

∀ψ ∈ H : lim U (t)ψ = U (t0 )ψ.
t→t0

Definition. Let U (·) : R → L(H) be a unitary, Abelian, one-parameter group (UAOPG).

Then its ( generator) is the linear map
Stone
A : DA →H
i
ψ 7→ Aψ := lim U (ε)ψ − ψ ,
ε→0 ε

where
i
Stone
:= ψ ∈ H | lim U (ε)ψ − ψ exists .

DA
ε→0 ε

Remark 12.7 . Note

i U (0 + ε)ψ − U (0)ψ 0
lim U (ε)ψ − ψ = i lim := i U (·)ψ (0),
ε→0 ε ε→0 ε
and so we can rewrite
0
S Stone
:= ψ ∈ H | i U (·)ψ (0) exists .

DA ≡ DA
0
Note also that U (·)ψ : R → H.

– 122 –
12.2 Stone’s Theorem
Theorem 12.8 (Stone’s Theorem). Let U (·) be a UAOPG that is strongly continuous and
whose group operation is the composition of maps, i.e.

U (t) ◦ U (s) = U (t + s)
U (0) = idH .

S → H is self adjoint on D S , and

Then its generator A : DA A

U (t) = exp(−itA).

Before proving this, consider the following.

Corollary 12.9. Given U (t) = exp(−itA) for some self adjoint A then the Spectral Theo-
rem tells us that U (t) is a UAOPG.

Proof. (i) First show U (t) ◦ U (s) = U (t + s):

U (t) ◦ U (s) := exp(−itA) ◦ exp(−isA)

Z
:= e−itλ e−isλ P (dλ)
ZR

= e−i(t+s) P (dλ)
R
=: U (t + s),

where Example 11.3 has been used.

We also have
U (0) := exp(0) = idH .
Now show Abelian property:
Z ∗
∗ −itλ
U (t) := e P (dλ)
Z R

= e+itλ P (dλ)
R
=: U (−t).

Then from the above we have U ∗ (t)U (t) = U (−t + t) = U (0) = idH . Then noticing that
kU (t)k = kU ∗ (t)k it follows that
p √
kU (t)k = k idH k = 1 = 1,

where we have used the fact that the norm is strictly positive to remove the negative root,
and so it is unitary.
Finally show that it is a group. This is easily done, and we have e = idH = U (0) and
[U (t)]−1 = U (−t).

– 123 –
S for some generator A and t ∈ R. Then
Corollary 12.10. Let ψ ∈ DA

[U (·)ψ]0 (t) = −iAU (t)ψ,

and
S S
U (t)DA = DA .

Try do proof later.

Now we can proceed with the proof of Stone’s Theorem. To do so, we will need to
show:

(i) The generator A is densely defined (otherwise we A∗ wouldnt be defined properly).

(ii) A is symmetric on DA
S and that it is essentially self adjoint.

(iii) U (t) = exp(−itA∗∗ ), from which it follows that A = A∗∗ and so it is self adjoint (as
A∗∗ is).

Proof. (Stone’s Theorem)

(i) Let ψ ∈ H. If we can show that an arbitrarily small neighbourhood N around ψ

S , then we know D S is dense in H. Consider the real family, for all
contains a ϕ ∈ DA A
τ ∈R Z τ
ψτ := drU (r)ψ,
0
which satisfies
ψτ
lim = ψ.
τ →0 τ
This is just the idea of the points in a neighbourhood, and so we know, therefore,
that there exists a τ0 such that ψτ0 ∈ N . Now consider
Z τ Z τ

U (ε)ψτ − ψτ := U (ε) drU (r)ψ − U (r)ψ
Z τ 0 Z 0τ
= drU (ε + τ )ψ − drU (r)ψ
0 0
Z ε+τ Z τ
= drU (r)ψ − drU (r)ψ
0 0
Z τ +ε Z τ Z τ
= drU (r)ψ + drU (r)ψ − drU (r)ψ
Zτ τ +ε Zε ε Z 0τ
= drU (r)ψ − drU (r)ψ − drU (r)ψ
τ τ 0
Z τ +ε Z ε
= drU (r)ψ − drU (r)ψ
τ Z ε 0 Z ε
= U (τ ) drU (r)ψ − drU (r)ψ
0 0
= U (τ ) − idH ψε .

– 124 –
Now using the fact that kU (τ )k = k idH k = 1 and so they’re bounded, we can take a
limit and push it through the operators. Thus we have
0 i
i U (·)ψτ (0) ≡ lim U (ε)ψτ − ψτ
ε→0 ε
1
= i U (τ ) − idH lim ψε
ε→0 ε

= i U (τ ) − idH ψ,

which is an element of H, and so we know that ψτ ∈ DA

S and therefore there exists a

ψτ0 ∈ N ∩ DA S.

(ii) Let ϕ, ψ ∈ DA
S , then

i
hϕ, Aψi := ϕ, lim U (ε)ψ − ψ
ε→0 ε

−i ∗
= lim U (ε) − idH ϕ, ψ
ε→0 ε

i
= lim U (−ε) − idH ϕ, ψ
ε→0 −ε

i
= lim U (ε) − idH ϕ, ψ
ε→0 ε
=: hAϕ, ψi,

where we have used the continuity of the inner product to move the limit in and out,
the result U ∗ (t) = U (−t), the fact that the identity is self adjoint and the fact that
we’re taking the limit to ‘ignore’ the minus signs on second to last line.
We now want to show that it is essentially self adjoint. Recalling Theorem 7.26, we
need to check if: for z ∈ C \ R that

ker(A∗ − z) = {0H } = ker(A∗ − z).

Let ϕ ∈ ker(A∗ − z) ∩ DA
S . Then for all ψ ∈ D S
∗ A
0
hϕ, U (·)ψi (t) = hϕ, [U (·)ψ]0 (t)i

= hϕ, −iAU (t)ψi

= −ihA∗ ϕ, U (t)ψi
= −ihzϕ, U (t)ψi
= −izhϕ, U (·)ψi(t)
=⇒ hϕ, U (·)ψi(t) = hϕ, ψie−izt ,

where we have used Corollary 12.10 and the fact that U (0) = idH . But, since z
is purely imaginary, the exponential is unbounded and so the RHS is unbounded.
However, the LHS is bounded (as U (·) is bounded) and so the only way the equality
holds is if hϕ, ψi = 0. Finally since we took all ψ ∈ H it follows that ϕ = {0H } and
so ker(A∗ − z) = {0H }. The proof for ker(A∗ − z) follows trivially from this result —
i.e. the RHS just becomes unbounded in the opposite direction.

– 125 –
(iii) We know that A is essentially self adjoint, which means that A∗∗ is self adjoint. Now
construct Z
∗∗
U (t) := exp(−itA ) =
e eitλ PA∗∗ (dλ).
R

Now let ψ ∈ S
DA ⊆ DA∗∗ and consider the real family

ψ(t) := exp(−itA∗∗ ) − U (t) ψ.

Then
ψ 0 (t) = − tA∗∗ exp(−itA∗∗ ) + iAU (t) ψ = −iA∗∗ ψ(t),

where we have used the fact that A = A∗∗ on DA

S . Then we have

(kψ(t)k2 )0 := hψ(t), ψ(t)i0

= 2 Rehψ(t), ψ 0 (t)i
= 2 Re − ihψ(t), A∗∗ ψ(t)i

= 0,

where we have used the fact that hψ, A∗∗ ψ(t)i ∈ R as A∗∗ is self adjoint. So we have
that kψ(t)k is a constant w.r.t. t. From the definition, we have ψ(0) = 0 and so
kψ(t)k = kψ(0)k = 0, which from the definition of the norm tells us ψ(t) = 0 for all
t. Finally it follows that

exp(−itA) =: U (t) = exp(−itA∗∗ ) =⇒ A = A∗∗ .

12.3 Domains of Essential Self Adjointness ("Cores")

Stone’s Theorem showed us that the generator A : DA S → H is self adjoint. Sometimes a

compromise in choosing the domain is in order, as we shall see in the two section’s time.

Corollary 12.11. Inspection of the part (ii) of the proof shows that if one considers A : D →
H for some dense D ⊆ DA S that also satisfies U (t)D = D , then we A is essentially self
A
adjoint on D.

12.4 Position, Momentum and Angular Momentum

Employ Stone’s Theorem to properly and easily define these three operators in quantum
mechanical systems. For the rest of this lecture we shall take H = L2 (R3 , λ) =: L2 where λ
is the Lebesgue measure.

Definition. The position operators, denoted Qj for j = 1, 2, 3, are defined as the generators
of
U j (·) : L2 → L2 ,
with
j
U j (t)ψ)(x) := ψ(x)e−itx ,

– 126 –
for (x) := (x1 , x2 , x3 ). That is, they are the self adjoint
Qj : DQ
S
j → L
2

with
(Qj ψ)(x) = xj ψ(x).
Remark 12.12 . Note clearly U (t) ◦ U (s) = U (t + s), U (0) = idL2 and kU (t)ψk = kψk which
tells us kU (t)k = 1, all of which are required for U (t) to be a UAOPG.
Definition. The momentum operators, denoted Pj for j = 1, 2, 3, are the generators of
Uj (·) : L2 → L2 ,
with
Uj (a)ψ)(x) := ψ(..., xj − a, ...),
i.e. they shift the j th slot to the right by a. That is they are the self adjoint operators on
their Stone domain that satisfy
Pj ψ = −i∂j ψ.
Remark 12.13 . Note this is exactly the definition we used for the action of the operator in
Lecture 9.
Definition. The orbital angular momentum operators, denoted Lj for j = 1, 2, 3, are the
generators of
Uj (·) : L2 → L2 ,
with
Uj (α)ψ)(x) := ψ(Dj (α)x),
where Dj (α) : R3 → R3 is the operator that describes the rotation about the j th axis by
angle α. They satisfy
(L1 ψ)(x) = −i(x2 ∂3 ψ − x3 ∂2 ψ)
(L2 ψ)(x) = −i(x3 ∂1 ψ − x1 ∂3 ψ)
(L3 ψ)(x) = −i(x1 ∂2 ψ − x2 ∂1 ψ)
Corollary 12.14. The spectrum for the orbital angular momentum is contained within the
integers; σ(Lj ) ⊆ Z for j = 1, 2, 3.
Proof. From Stone’s theorem we have
Uj (α) = exp(−iαLj ),
which together with Dj (α + 2π) = Dj (α) gives
exp(−i2πLj ) = idH .
Then, from the fact that Lj is self adjoint, we can use the Spectral theorem to decompose
both sides Z Z
−i2πλ
e PLj (dλ) = PLj (dλ),
R R
and so λ ∈ Z.

– 127 –
12.5 Schwartz Space S(Rd )
As we have just seen, Stone’s theorem gives us a nice way to define the position, momen-
tum and orbital angular momentum operators. However there are two problems with the
definitions we have, both of which relate to their Stone domains. They are
(i) DQ
S 6= D S 6= D S 6= D S
P L Q

(ii) DQ Q◦Q 6= DQ◦Q◦Q 6= ... and similarly for P and L.

S 6= D S S

This, at first, might not seem like such a big deal but on a second look we see that
it means havoc when it comes to trying to define the QM version of kinetic energy as
(P ◦ P )/2m. The problem is especially bad when it comes to considering commutators, as
highlighted before.
We get around this problem using the compromise given in Corollary 12.11.
Definition. The Schwartz Space on Rd , denoted S(Rd ), is the vector space with set

S(Rd ) := {ψ ∈ C ∞ (Rd ) | sup |xα (∂β ψ)(x)| < ∞, ∀α, β ∈ N×d

0 },
x∈Rd

where
N×d N0 × ... × N0 ,
0 := | {z }
d-fold
and

xα ≡ x(α1 ,...,αd ) := (x1 )α1 ...(xd )αd ,

d
∂β ≡ ∂(β1 ,...,βd ) := (∂1 )β1 ...(∂d )β .

Remark 12.15 . The Schwartz Space is also known as the space of rapidly decaying test
functions.
Remark 12.16 . Clearly the space Cc∞ (Rd ), as defined in footnote 15 in Lecture 9, is a
contained within S(Rd ).
Lemma 12.17. The Schwartz space is closed under pointwise multiplication; if ψ, ϕ ∈
S(Rd ) then ψ • ϕ ∈ S(Rd ). In fact we have the Schwartz algebra (S(Rd ), +, ·, •).
Proof. This result follows simply from the so called Leibniz Rule, which is an extension of
the product rule.22

sup xα ∂β (ψ • ϕ) (x) = sup xα ∂β (ψ) • ϕ + ψ • ∂β (ϕ) (x)

x∈Rd x∈Rd
≤ sup xα ∂β (ψ) • ϕ (x) + sup xα ψ • ∂β (ϕ) (x)

x∈Rd x∈Rd
< ∞.

Then, using the fact that the pointwise multiplication of two smooth functions is smooth,
we have ψ • ϕ ∈ S(Rd ). Finally using the linearity of everything involved we get the
algebra.
22
See Dr. Schuller’s Lecture’s on the Geometric Anatomy of Theoretical Phsyics for a definition in
context.

– 128 –
Lemma 12.18. For 1 ≤ p ≤ ∞ we have S(Rd ) ⊂ Lp (Rd ).

Proof. Let ψ ∈ S(Rd ). Then |ψ(x)| < ∞ for all x ∈ Rd , and so it is integrable. Then
Corollary 5.19 tells us that it is measurable with respect to the Borel σ-algebras, so ψ ∈
L1 (Rd ). Then finally from L1 (Rd ) ⊂ Lp (Rd ) for all p > 1, the result follows.

Lemma 12.19. One can show that the Fourier Transform is a linear isomorphism from
S(Rd ) onto itself.23

Theorem 12.20. The Schwartz space as defined above satisfies

(i) S(Rd ) ⊆ L2 (Rd , λ) is dense,

(ii) S(R3 ) ⊆ DQ
S , D S , D S is dense,
P L

(iii) Qj : S(R3 ) → S(R3 ) is essentially self adjoint. Same for Pj and Lj .

Remark 12.21 . From the last condition we see that we can repeatedly apply the operators,
in any order, to a system. This fixes our problem above.

23
See lecture 18.

– 129 –
13 Spin

In the previous lecture we defined orbital angular momentum. The emphasis on ‘orbital’ was
not a mistake; this lecture aims to discuss what is referred to as general angular momentum
(or just angular momentum) in QM. This latter case gets its name from the fact that any
concrete set of three operators, {J1 , J2 , J3 } say, obey analogous commutation relations to
{L1 , L2 , L3 }.
Recall, for ψ ∈ S(R3 ) we have

Li : S(R3 ) → S(R3 ),

essentially self adjoint with

(L1 ψ)(x) = −i(x2 ∂3 ψ − x3 ∂2 ψ),

(L2 ψ)(x) = −i(x3 ∂1 ψ − x1 ∂3 ψ),
(L3 ψ)(x) = −i(x1 ∂2 ψ − x2 ∂1 ψ).

We can, therefore, calculate their commutation relations.

Lemma 13.1. The orbital angular momentum operators obey the following commutation
relations:
[L1 , L2 ] = iL3 ,
[L2 , L3 ] = iL1 ,
[L3 , L1 ] = iL2 .

Proof. The proof follows from direct computation. Consider

[L1 , L2 ]ψ = (L1 ◦ L2 − L2 ◦ L1 )ψ
= (−i)2 (x2 ∂3 − x3 ∂2 )(x3 ∂1 ψ − x1 ∂3 ψ) − (−i)2 (x3 ∂1 − x1 ∂3 )(x2 ∂3 ψ − x3 ∂2 ψ)
= − x2 ∂1 ψ + x2 x3 ∂3 ∂1 ψ − x2 x1 ∂32 ψ − (x3 )2 ∂2 ∂1 ψ + x3 x1 ∂2∂3 ψ

+ x2 x3 ∂1 ∂3 ψ − (x3 )2 ∂1 ∂2 ψ − x1 x2 ∂32 ψ + x1 x3 ∂3 ∂2 ψ + x1 ∂2 ψ

= x1 ∂2 ψ − x2 ∂1 ψ
= iL3 ψ,

where we have used the fact that S(R3 ) ⊂ C ∞ (R3 ) ⊂ C 2 (R3 ), and so we can swap derivative
order, i.e.
∂1 ∂2 ψ = ∂2 ∂1 ψ.
The same method is used for the other two commutation relations.

Remark 13.2 . The vector space with set V := spanC {L1 , L2 , L3 } can be defined, and we
have that iLj ∈ V , and so the above tells us that (V, +, ·, [·, ·]) is the orbital angular
momentum Lie algebra.

– 130 –
Remark 13.3 . We can re-write the commutation relations in the compact and convenient
form
[Li , Lj ] = iijk Lk ,
where ijk is the totally antisymmetric Levi-Civita symbol, defined as

+1 if (i, j, k) is an even permutation of (1, 2, 3)


ijk = −1 if (i, j, k) is an odd permutation of (1, 2, 3)

otherwise.

0

Proposition 13.4. It is not possible have a common eigenvector between the operators
(L1 , L2 , L3 ).

Proof. Let V := spanC {L1 , L2 , L3 }. Now assume that ψ ∈ D \ {0H } is an eigenvalue of

both L1 and L2 with eigenvalues λ, µ ∈ C, respectively. Then we have

[L1 , L2 ]ψ := L1 (L2 ψ) − L2 (L1 ψ)

= µL1 (ψ) − λL2 ψ
= µλψ − λµψ
= 0,

where we have used the linearity of the operators. It follows from the commutation relations
that L3 ψ = 0. However, from the other commutation relations we then have

[L2 , L3 ]ψ = iL1 ψ
L2 (L3 ψ) − L3 (L2 ψ) = iλψ
L2 (0) − µL3 ψ = iλψ
0 − µ0 = iλψ
=⇒ λ = 0,

and similarly you can show µ = 0. It follows then that for any D ∈ V we have Dψ = 0,
which can only be true is ψ = 0H . But this contradicts the opening assumption and so it
can’t be true.

Remark 13.5 . People often say that "two non-commuting operators have no common eigen-
vectors", however this statement is not strictly true. What is meant is "two operators, whose
commutator does not contain the zero vector in its range, do not have common eigenvec-
tors." This is subtly different, however the distinction is important. For example, in the
previous proposition if we instead had [L1 , L2 ] = iL3 and [L2 , L3 ] = 0 = [L3 , L1 ], we would
not need to require λ = 0 = µ, and so, unless further constraints were placed on the system
of operators, it is possible that ψ is a common eigenvector to L1 and L2 .

– 131 –
13.1 General Spin
At this point we might ask why we are bothering to work out these commutation relations?
After all they appear to be of no use when it comes to calculating things such as the spectra
of the operators (as is evident by Corollary 12.14). The answer to this is that we want to see
what information we can obtain about the system (specifically its spectrum) using solely the
commutation relations, as then any other set of observables that shares these commutation
relations immediately obey the same results.
To emphasise, given a Lie algebra that contains three operators, S1 , S2 , S3 say, with

Sj : D → D,

for some D ⊆ H, that obey commutation relations analogous to those of Lemma 13.1, will
instantly satisfy any results we derive, using only the commutation relations, for the orbital
angular momentum.

Example 13.6 . An example of such a set of operators are the so-called Pauli spin algebra,
which has
1
Si := σi ,
2
with D = H = C2 , where
! ! !
01 0 −i 1 0
σ1 = , σ2 = , σ3 = ,
10 i 0 0 −1

known as the Pauli spin matrices. It is easily checked, through the rules of matrix multi-
plication, that this algebra obeys the correct commutation relations. This is an example of
a so-called spin- 12 system.

Remark 13.7 . We can not expect the commutation relations to necessarily tell us everything
about the spectrum (or any other quantity we try to calculate) as they can be derived
by several different operator sets with potentially differing spectra. But, as said above,
whatever we can infer from the commutation relations alone must hold for all the operator
sets.

13.2 Derivation of Pure Point Spectrum

We start from a general Lie algebra with our required conditions. We shall denote the
operators by J1 , J2 , J3 , however they need not be the orbital angular momentum operators.
Equally the domain, D, is left arbitrary, up the condition that the operators are at least
essentially self adjoint on them.

Definition. A Casimir operator for the algebra (V, +, ·, [·, ·]), is a symmetric operator

Ω: D → D

that commutes with every element in V .

– 132 –
Remark 13.8 . Note, due to the bilinearity of the commutator, we only need to check that
Ω commutes with the basis elements of V .

Proposition 13.9. The operator

Ω := J1 ◦ J1 + J2 ◦ J2 + J3 ◦ J3

is a Casimir operator for the algebras we’re considering.

Proof. The symmetric part follows trivially from the fact that Corollary 12.11 tells us
J1 , J2 , J3 are all symmetric. Next let ψ ∈ D and consider

[Ω, J1 ]ψ := [J1 ◦ J1 + J2 ◦ J2 + J3 ◦ J3 , J1 ]ψ
= [J1 ◦ J1 , J1 ]ψ + [J2 ◦ J2 , J1 ]ψ + [J3 ◦ J3 , J1 ]ψ
= (J1 ◦ J1 ◦ J1 )ψ − (J1 ◦ J1 ◦ J1 )ψ
+(J2 ◦ J2 ◦ J1 )ψ − (J1 ◦ J2 ◦ J2 )ψ
+(J3 ◦ J3 ◦ J1 )ψ − (J1 ◦ J3 ◦ J3 )ψ
= (J2 ◦ J1 ◦ J2 )ψ + (J2 ◦ [J2 , J1 ])ψ − (J1 ◦ J2 ◦ J2 )ψ
+(J3 ◦ J1 ◦ J3 )ψ + (J3 ◦ [J3 , J1 ])ψ − (J1 ◦ J3 ◦ J3 )ψ
= (J1 ◦ J2 ◦ J2 )ψ + ([J2 , J1 ] ◦ J2 )ψ + (J2 ◦ [J2 , J1 ])ψ − (J1 ◦ J2 ◦ J2 )ψ
+(J1 ◦ J3 ◦ J3 )ψ + ([J3 , J1 ] ◦ J3 )ψ + (J3 ◦ [J3 , J1 ])ψ − (J1 ◦ J3 ◦ J3 )ψ
= −i(J3 ◦ J2 )ψ − i(J2 ◦ J3 )ψ + i(J2 ◦ J3 )ψ + i(J3 ◦ J2 )ψ
= 0,

which because ψ ∈ D was arbitrary tells us [Ω, J1 ] = 0. The same method gives [Ω, J2 ] =
0 = [Ω, J3 ].

Definition. Let J1 , J2 , J3 be three operators that satisfy our conditions. Then define

J+ := J1 + iJ2 ,
J− := J1 − iJ2 ,

known as the ladder operators, for a reason that will soon become clear.

Remark 13.10 . We can choose to consider the set {J+ , J− , J3 } in place of the set {J1 , J2 , J3 }
while still keeping all the information — as we can simply reconstruct J1 and J2 from J+
and J− . Note, however, in doing this we have broken the symmetry of the algebra (in the
sense that none of the Jj s are special, they all obey the same commutation relations) by
singling out J3 , while taking linear combinations of J1 and J2 . Indeed we did not need to
make this choice of symmetry breaking, but in fact we could have chosen to keep J1 while
defining J+ and J− as linear combinations of J2 and J3 . Importantly, the results that follow
will hold equally for whichever J we choose, and so in order to stick with convention we pick
J3 . Note also that we no longer have a set of observables as (J+ )∗ = J− and (J− )∗ = J+ .

– 133 –
Lemma 13.11. J+ and J− satisfy the following commutation relations

[J+ , J− ] = 2J3 ,

[J3 , J± ] = ±J± ,
[Ω, J± ] = 0.

Proof. From direct computation. Not done here to save space.

Lemma 13.12. We can rewrite our Casimir operator as

Ω = J+ ◦ J− + J3 ◦ (J3 − idD ),

or equally as
Ω = J− ◦ J+ + J3 ◦ (J3 + idD ).

Proof. We shall just show the first one:

J+ ◦ J− + J3 ◦ (J3 − idD ) := (J1 + iJ2 ) ◦ (J1 − iJ2 ) + J3 ◦ (J3 − idD )

= J1 ◦ J1 + J2 ◦ J2 − i(J1 ◦ J2 − J2 ◦ J1 ) + J3 ◦ J3 − J3
= J1 ◦ J1 + J2 ◦ J2 + J3 ◦ J3 − i[J1 , J2 ] − J3
= J1 ◦ J1 + J2 ◦ J2 + J3 ◦ J3 − (i)2 J3 − J 3
= Ω.

Remark 13.13 . Both of the expressions in the above definition are always true. It is not
that one is true under certain circumstances and then the other is true. This is an important
observation that we shall use.

At this point we might wonder why we are going through so much effort introducing the
Casimir, when what we’re looking for is the spectra of the operators J1 , J2 , J3 . The answer
is to make the problem seemingly more complicated by now considering only eigenvectors
that are common to both J3 and Ω. Note it is necessary that they commute if they are to
have common eigenvectors — as is easily verified from the definition of the commutator.
That is we want to find a ψλ,µ ∈ D \ {0}24 such that

J3 ψλ,µ = µψλ,µ
Ωψλ,µ = λψλ,µ ,

where the subscript is included in order to label the eigenvector by its eingenvalues.
Again this appears to be a more complicated problem — we now not only need to
check our ψ is an eigenvector of J3 but we also need to check that it’s an eigenvalue of Ω.
However, we can show that every eigenvector of J3 (and equally for J1 and J2 ) is also an
eigenvector of Ω. For a proof of this see Peter Ferguson’s well detailed answer on Quora.
24
Recall that an eigenvector can not be the zero-vector by definition. We shall use this later.

– 134 –
We might also ask whether this will give us the spectrum of J3 anyways, as J3 is only
essentially self adjoint on our Schwartz domain, D. We are OK, however as (J3 )∗∗ is self
adjoint on the D and so we could just use this instead, where the operation is defined to
be the same as J3 — just as we did for the momentum operator in Lecture 9.
Lemma 13.14. The eigenvalues for these common eigenvectors satisfy

λ ≥ |µ|(|µ| + 1).

Proof. We shall consider both cases for the rewriting of Ω, however, as explained above, the
cases bracket does not mean one is true under certain conditions and the other otherwise,
they are both true. We shall also drop the ◦ symbols to lighten notation. Thus we have

λhψλ,µ , ψλ,µ i = hψλ,µ , Ωψλ,µ i

(

ψλ,µ , J+ J− + J3 (J3 − idD ) ψλ,µ
=

ψλ,µ , J− J+ + J3 (J3 + idD ) ψλ,µ
(
hJ− ψλ,µ , J− ψλ,µ i + µ(µ − 1)hψλ,µ , ψλ,µ i
=
hJ+ ψλ,µ , J+ ψλ,µ i + µ(µ + 1)hψλ,µ , ψλ,µ i,

where we have used the fact that (J− )∗ = J+ and vice versa.25
Now recalling that ψλ,µ is an eigenvector, and so, by definition, not the zero vector we
know the inner product is positive definite, and thus we can divide by it, giving

 kJ− ψλ,µ k2 + µ(µ − 1)
2
λ = kJ+kψk ψ λ,µ k2
 2kψk
+ µ(µ + 1).

Finally, from the fact that the norm is non-negative definite (i.e. the first term in each case
is either positive or vanishes) we have
(
µ(µ − 1)
λ≥
µ(µ + 1)
(
−µ(−µ + 1)
=
µ(µ + 1),

and so it follows that λ ≥ |µ|(|µ| + 1).

Lemma 13.15. The elements J± ψλ,µ are common ‘eigenvectors’ of Ω and J3 with eigen-
values λ and (µ ± 1), respectively.
Proof. First consider Ω. We have

ΩJ± ψλ,µ = J± Ωψλ,µ + [Ω, J± ]ψλ,µ

= J± (λψλ,µ )
= λ(J± ψλ,µ ),
25
Really what you need to do is expand out J− and J+ in terms of J1 and J2 and then take the adjoint.
The result holds.

– 135 –
where we used the fact that [Ω, J± ] = 0, and the linearity of J± .
Now consider J3 ,

J3 J± ψλ,µ = J± J3 ψλ,µ + [J3 , J± ]ψλ,µ

= µ(J± ψλ,µ ) ± J± ψλ,µ
= (µ ± 1)(J± ψλ,µ ),

where again we used the commutator, [J3 , J± ] = ±J± .

Remark 13.16 . The previous allows us to conclude that J± ψλ,µ ∝ ψλ,µ±1 .

Remark 13.17 . This is why J± are known as ladder operators, with J+ known as the raising
operator and J− the lowering operator. As we see these names derive from the eigenvalues
they produce as eigenvectors of Jz .
If the Jz eigenvalue of ψλ,µ corresponds to the µ-th rung of a ladder then the eigenvalue
of J+ ψλµ corresponds to the (µ + 1)-th rung, and J− ψλ,µ the (µ − 1)-th. Note each one of
the rungs is separated by exactly the same distance, and that we get the (µ + n)-th rung
from (J+ )n ψλ,µ and similarly for the (µ − n)-th rung.

J+ ψλ,µ (µ + 1)
Eigenvectors

Eigenvalues

ψλ,µ µ

J− ψλ,µ (µ − 1)

The next question would be is this a ‘proper’ ladder; that is does it have a top and bottom
rung or does it continue forever? The answer comes in the form of the next lemma.

Lemma 13.18. There exists a ψλ,µ such that J+ ψλ,µ = 0. Equally there exists a ψλ,µ such
that J− ψλ,µ = 0.

Proof. We know from Lemma 13.14 that |µ|(|µ| + 1) ≤ λ holds for any common eigenvector
of Ω and J3 . We see from Lemma 13.15 that (J± )n ψλ,µ is such a common eigenvector, and
so must obey |µ ± n|(|µ ± n| + 1) ≤ λ. However, λ is unchanged by this repeated application
of the ladder operators, and so, unless remedied, this inequality will eventually be broken
– that is we need to somehow cap the available n values.
Consider first the raising operator. In this case µ + n gets bigger and bigger, and so
we need to cap n from above. In other words, we require there to be an m ∈ N such that
for all n > m, (J+ )n ψλ,µ = 0. This fixes our problem as this corresponds to the zero vector

– 136 –
and so, by definition, it cannot be a eigenvector, and the λ inequality no longer need hold.
We define
ψλ,µ := (J+ )m ψλ,µ .
The idea is exactly the same is true for the lowering operator, however now µ − n is
getting smaller and smaller, and so its modulus (after n > µ is reached) gets bigger and
bigger. So again we need to cap n from above. We require there to be a ` ∈ N such that
for all n > `, (J− )n ψλ,µ = 0. We define

ψλ,µ := (J− )` ψλ,µ .

Note we do not have any a priori relation between the values of m and `. To use the
ladder analogy, m is the number of rungs above µ-th rung and ` is the number of rungs
below the µ-th rung. For a given λ, the highest value of µ is denoted µ(λ) and the lowest
value µ(λ).

Remark 13.19 . Note the above tells us that J± ψ are not strictly eigenvalues, as it could be
the zero vector. This is why we wrote ‘eigenvectors’ in inverted commas in Lemma 13.15.

Proposition 13.20. The maximum and minimum values of µ satisfy

(i) λ = µ(λ) µ(λ) + 1 ,

(ii) µ(λ) = −µ(λ),

N0
(iii) µ(λ) ∈ 2 .

Proof. (i) From the proof of Lemma 13.14, and the fact that J+ ψλ,µ(λ) = 0, and so
kJ+ ψλ,µ(λ) k = 0, we have

λ = µ(λ) µ(λ) + 1

(ii) Repeating the above argument but with the fact that kJ− ψλ,µ(λ) k = 0, we have

λ = µ(λ) µ(λ) − 1

= −µ(λ) − µ(λ) + 1 ,

and so µ(λ) = −µ.

(iii) From the previous, along with the fact that µ(λ) − µ(λ) ∈ N0 this result follows
trivially.

Remark 13.21 . In order to be consistent with the literature we shall introduce the following
relabelling
j := µ(λ), m := µ.
Note we have j ∈ 2 .
N0

– 137 –
Theorem 13.22. The common eigenvectors of Ω and J3 come as families ψj(j+1),m , where
m = −j, −j + 1, ..., j − 1, j. The eigenvalue j(j + 1) is associated to Ω and m is associated
to J3 .

One normally normalises these eigenvectors and defines

ψj(j+1),m
Φj,m := .
kψj(j+1),m k

Then, from Lemma 8.8 and the fact that the eigenvectors have distinct eigenvalues, we have

hΦj,m , Φk,n i = δjk δmn .

Z
Corollary 13.23. We have m ∈ 2.

Proof. This comes from just allowing j ∈ N0

2 to be any element and then using m =
−j, ..., j.

Proposition 13.24. The the common eigenvectors Φj,m satisfy

p
J± Φj,m = j(j + 1) − m(m ± 1)Φj,m±1 .

Proof. We shall show this for J+ , the method of J− follows analogously. Recall from
Lemma 13.12 that
J− J+ = Ω − J3 (J3 + idD ).
Now consider

hJ+ Φj,m , J+ Φj,m i = hΦj,m , J− J+ Φj,m i

= hΦj,m , ΩΦj,m i − hΦj,m , J3 (J3 + idD )Φj,m i
= j(j + 1)hΦj,m , Φj,m i − m(m + 1)hΦj,m , Φj,m i
= j(j + 1) − m(m + 1).

Combining this with Remark 13.16 allows us to conclude the result.

13.3 Pure Spin-j Systems

Definition. A quantum mechanical system is called a pure spin-j system if its Hilbert
space is (2j + 1)-dimensional that possesses an orthonormal eigenbasis {Φj,m } for the three
operators J1 , J2 , J3 defined on H.

Corollary 13.25. The Hilbert space is isomorphic to C2j+1 .

Corollary 13.26. For a pure spin-j system the spectrum of the operators is

σ(Ji ) = {−j, −j + 1, ..., j − 1, j},

for i = 1, 2, 3.

Example 13.27 .

– 138 –
j σ(Ji ), i = 1, 2, 3
0 0
1/2 {−1/2, 1/2}
1 {−1, 0, 1}

Remark 13.28 . When you introduce spin to a particle, its Hilbert space becomes a product
space. For example for an electron (spin-1/2) in R3 its Hilbert space is

He = L2 (R3 ) ⊗ C2 .

We shall return to this and expand on it in the next lecture.

Remark 13.29 . You can also have non-pure spin systems. For the orbital angular momen-
tum, you take a direct sum of the Hilbert spaces. That is if Hj is the Hilbert space associated
to the pure spin-j system then the composite system’s Hilbert space is
M
Hcomp = Hj .
j

We shall return to this in two lectures time.

– 139 –
14 Composite Systems
Recall Axiom 1, which says that to every quantum system there is an underlying Hilbert
space. The question we now want to ask is: Let H1 be the Hilbert space associated to one
system and H2 be the Hilbert space associated to another. What is the underlying Hilbert
space associated to the composite system?
To clarify what we mean, imagine having a proton and an electron. We first look at the
proton by itself and call this system one. We then look at the electron separately and call
that system two. We now want to look at both of them together, but we wish to use the
fact that we have already studied them separately to simplify the problem. It may seem
‘natural’ to model the composite H as the so called direct sum, which as a set is26
H1 ⊕ H2 := {(ψ, ϕ) | ψ ∈ H1 , ϕ ∈ H2 },
and where the linearity is inherited from H1 and H2 , namely
(aψ1 + ψ2 , bϕ1 + ϕ2 ) = ab(ψ1 , ϕ1 ) + a(ψ1 , ϕ2 ) + b(ψ2 , ϕ1 ) + (ψ2 , ϕ2 ).
This is what we do in classical systems and it tells us that if we know everything about the
states27 of our two systems, then we also know everything about the states of the composite
system.
However, as with all things quantum, things are more complicated, and the above is
not the case. The main problem comes from the fact that not all linear combinations of
elements of the form (ψ, ϕ) can also be written in that form.
Example 14.1 . Let ψ1 , ψ2 ∈ H1 , ϕ1 , ϕ2 ∈ H2 and a, b ∈ C. Then, assuming the linearity as
above, we have
a(ψ1 , ϕ1 ) + b(ψ2 , ϕ2 ) = (aψ1 , ϕ1 ) + (ψ2 , bϕ2 )
= (aψ1 + ψ2 , ϕ1 + bϕ2 )
= a(ψ1 , ϕ1 ) + b(ψ2 , ϕ2 ) + ab(ψ1 , ϕ2 ) + (ψ2 , ϕ1 ),
a clear problem.
Note this example actually tells us that we H1 ⊕ H2 is not closed under the linearity,
and so would not be a vector space. We could just restrict ourselves to elements that
do obey these rules, however, as we shall see when considering entanglement, we require
elements of this form in our underlying Hilbert space.
This calls for a slight refinement of axiom one. We add the addendum28 :

If a quantum system is composed of two (and hence, by induction, any finite number)
of ‘sub’systems, then its underlying Hilbert space is the tensor product space H1 ⊗H2 ,
equipped with a inner product.

26
This definition holds as we are only taking the direct product of two spaces, and so the index set is
finite. See wiki for details on this.
27
Recall that the elements of the Hilbert space are not the states, but are associated to them. We shall
return to this at the end of the lecture.
28
We shall define what these new terms are in the next section.

– 140 –
Example 14.2 . Remark 13.28 is an example of such a composite system.

14.1 Tensor Product of Hilbert Spaces

In order to give a nice definition for the tensor product of two vector spaces, we first need
to introduce the so called free vector space.

Definition. Let V be a F-vector space and let B ⊆ V be a generating subset of V (i.e.

any element of V can be obtained via finite linear combinations of elements of B). The the
free vector space is
F (B) := spanF (B),
i.e. the set of all linear combinations of elements of B.

Lemma 14.3. Every vector space is a free vector space with B being a Hamel basis.

Remark 14.4 . Note it need not be true that F (B) = V , as it might be the case that the
same element in V is reached via two different linear combinations of elements of B. In
fact if F (B) = V , then B is just a Hamel basis.

The free vector space for vector spaces might seem almost redundant, given that every
vector space has a basis. However if your vector space is countably infinite then such a
basis might be incredibly difficult to construct. However you can simply take the entire set
for B and construct the free vector space F (V ), which will be a huge set, mind. Note, then,
that any linear combination of elements in this set is automatically still in the set, and so
it is indeed a vector space.

Definition. Let V and W be two F-vector spaces, and let A ⊆ V and B ⊆ W be generating
subsets. Then we define their tensor product as the vector space with set

V ⊗ W := F (A × B)/∼ ,

where × is the Cartesian prodcut and ∼ is an equivalence relation such that: if a, a1 , a2 ∈ A,

b, b1 , b2 ∈ B and f ∈ F then

(i) (a, b) ∼ (a, b),

(ii) (a1 , b1 ) + (a1 , b2 ) ∼ (a1 , b1 + b2 ) and (a1 , b1 ) + (a2 , b1 ) ∼ (a1 + a2 , b1 ), and continued
by induction,

(iii) f (a, b) ∼ (f a, b) and f (a, b) ∼ (a, f b).

(iv) Combinations of (ii) and (iii), e.g. (a1 , b1 ) + f (a1 , b2 ) ∼ (a1 , b1 + f b2 ).

Remark 14.5 . Note the equivalence relation looks a lot like a linearity condition on V ⊗ W ,
however on closer inspection it is not quite. The linearity condition that make V ⊗ W into
a vector space is simply
f (a1 , b1 ) + (a2 , b2 ) ∈ V ⊗ W.

– 141 –
This, in itself, does not need to satisfy the equivalence relation. However, if we did not
include it we could end up with a huge redundancy in elements, as a repercussion of Re-
mark 14.4. This equivalence relation makes the corresponding set of equivalence classes a
vector space in the way we normally think of them (there is no repeated elements).

This is exactly the type of structure we need to overcome the problem highlighted
before (that not all linear combinations can be expressed as a single term), as now we only
require that linear combinations of linear combinations are linear combinations, which they
obviously are.

Proposition 14.6. Let H1 and H2 be our two vector spaces. We can define the map

+H1 ⊗H2 : (H1 ⊗ H2 ) × (H1 ⊗ H2 ) → (H1 ⊗ H2 )

[(ψ, ϕ1 )], [(ψ, ϕ2 )] 7→ [(ψ, ϕ1 )] +H1 ⊗H2 [(ψ, ϕ2 )] := [(ψ, ϕ1 ) + (ψ, ϕ2 )],

where the additions inside the brackets are w.r.t. H1 and H2 .

Proof. We need to show this is well defined. We shall write +12 now in order to lighten
notation. Consider it case by case.